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CROSS REFERENCES TO RELATED APPLICATIONS This application is a divisional of U.S. application Ser. No. 09/110,503, filed Jul. 6, 1998, which claimed priority based on U.S. provisional application No. 60/052,587, filed Jul. 15, 1997, the disclosures of both said applications being incorporated by reference herein. STATEMENT AS TO RIGHTS TO INVENTIONS MADE UNDER FEDERALLY SPONSORED RESEARCH AND DEVELOPMENT Not applicable. BACKGROUND OF THE INVENTION The present invention is in the field of biotechnology and specifically relates to the study of the physiological and physiochemical processes which govern and underlie the formation, growth and resorption of human and animal bone. In particular the invention provides novel means for the study of responses of the mammalian musculoskeletal system to stress and potentially may lead to the discovery of novel substances produced by bone during these responses. The instant system may lead to a better understanding of diseases such as osteoporosis and the perfusion chamber means provides means for the study of the effects of drugs and other substances added to the perfused medium. It has been known for over 150 years that bone responds to mechanical loading. Although the effects of exercise and mechanical loading on the musculoskeletal systems have been well documented, the actual mechanisms by which mechanical loading acts at the cellular level in the maintenance of skeletal integrity are not completely understood. Although greater attention is being given to exercise and nutrition as a means of preventing and/or treating osteoporosis, the regulatory mechanisms that control skeletal response to mechanical loading, growth factors and nutrition are not yet delineated. There is speculation about the biophysical structure and properties of the sensory and biochemical and molecular biological mechanism of mechano-transduction. When controlled loads of a given magnitude and frequency are applied, in vivo, either in an isolated wing preparation or a rat tibia, bone mineral density is known to increase to an extent which is approximately proportional to the load applied. However, according to the prior art, it is not possible to assess quantitatively the bone-specific regulatory control product and their mechanisms nor to monitor the bone production of local growth factors and cytokines, in these in vivo preparations. Whilst cell culture preparations do permit an investigator to quantify second messengers, cytokines and local growth factors, they do not permit one to monitor the responses of bone cells to mechanical deformation of the bone matrix which are so important in maintaining and/or remodeling of the skeletal system. Although growth factors have been shown to enhance the development of new bone, clearly and without the presence of mechanical loading, under these circumstances, the new matrix is not formed along lines of strain and it is that feature, in life, which induces maximum integrity of the new bone so formed. The present authors have been associated with previous work in which the viability of osteoblasts from 2 to 4 week old pigs was successfully maintained, in culture, for 68 days. Careful consideration of these findings led to the hypothesis that, in a suitable novel system, which would permit continuous perfusion and mechanical loading of suitable explanted samples of trabecular bone from mature pigs, viability might be maintained for 10 to 12 days or longer. If this were to be achieved, such a time frame would permit measurements of the rate of bone formation and resorption of the trabecular bone, not available using the systems, apparatus and methods of the prior art. Further, such a novel system would be applicable to the study of human bone. Up to now, prior art apparatus and systems for investigating bone have either comprised cell culture apparatus of a variety of well-known types or mechanical means for applying three point and four point bending forces to a biological test subject. An example of the three point type is disclosed in U.S. Pat. No. 5,406,853 to Lintilhac and Vesecky and an example of the four point type is disclosed in U.S. Pat. No. 5,383,474 to Recker and Akhter. The present authors are not aware of any prior art system or apparatus which provides means for simultaneous, contemporaneous and continuous study of axially loaded viable mammalian bone undergoing concurrent continuous perfusion and the effluent medium therefrom. References 1. Baron R., Vignery A., Neff L., Silverglate A., Santa Maria A. (1983). Processing of undecalcified bone specimens for bone histomorphometry. In: ed, Recker R. R., Bone Histomorphometry: Techniques and Interpretation CRC, Boca Raton, Fla., 13-35. 2. Brighton C. T., Sennett B. J., Farmer J. C., Ianotti J. P., Hansen C. A., Williams J. L., Williamson J. The inositol phosphate pathway as a mediator in the proliferative response of rat calvarial bone cells to cyclical biaxial mechanical strain. J. Orthop. Res. 10:385-393; 1992. 3. Carvalho R. S., Scott J. E., Suga D. M., Yen E. H. K. Stimulation of signal transduction pathways in osteoblasts by mechanical strain potentiated by parathyroid hormone. J. Bone Min. Res. 9::999-1011; 1994. 4. Crenshaw T. D., Thomson B. M., Noble B. S., Milne J. S. and Loveridge N. Prostaglandin E2 inhibits proliferation of porcine progenitor osteoblast cells. J. Bone Min. Res. 8(1):S362, 1993. 5. Currey J. The Mechanical Adaptation of bones. Princeton N.J.: Princeton University Press, 1984. 6. Dalsky G. P., Stocke K. S., Ehsani A. A., Slatopolsky E., Lee W. C. and Birge S. J. Weight-bearing exercise training and lumbar bone mineral content in post menopausal women. Ann. Intern. Med. 108:824-828, 1988. 7. E1 Haj A. J., Minter S. L., Rawlinson S. C. F., Suswillo R. and Lanyon L. E. Cellular responses to mechanical loading into vitro. J. Bone Min. Res. 5:923-932, 1990. 8. Frost H. M. (1983). Bone histomorphometry: analysis of trabecular bone dynamics. In: ed, Recker R. R. Bone Histomorphometry: Techniques and Interpretation CRC Press, Boca Raton, Fla., 109-142. 9. Gleeson P. B., Protas E. J., Le Blanc A. D., Schneider V. S. and Evans H. J. Effects of weight lifting on bone mineral density in premenopausal women. J. Bone, Min. Res. 5:153-157, 1990. 10. Hock J. M., Centrella M. and Canalis E. Insulin-like growth factor I (aIGF-1) has independent effects on bone matrix formation and cell replication. Endocrin. 122:254-260; 1988. 11. Hsieh H-J., Li N. Q., Frangos J. A. Shear stress increases endothelial platelet-derived growth factor mRNA levels. Am. Physiol J. 260:H642-646; 1991. 12. Jones D. B., Nolte H., Scholubbers J. G., Turner E., Veltel D. Biochemical signal transduction of mechanical strain in osteoblast-like cells. Biomaterials. 12:101-110; 1991. 13. Lanyon L. E. Control of bone architecture by functional load bearing. J. Bone Min. Res. 7:S369-S375, 1992. 14. Murray D. W. and Rushton N. The effect of strain on bone cell prostaglandin E2 release: a new experimental method. Calcif. Tissue Int. 47:35-39, 1990. 15. Mundy G. R. Bone resorbing cells. In: Primer on the Metabolic Bone Diseases and Disorders of Mineral Metabolism. Favus M. J. (ed), Kelseyville, Calif. American Society for Bone and Mineral Research (PP.18-22) 1990. 16. Notelovitz M., Martin D., Tesar R., Khan F. Y., Probart C., Fields C. and McKenzie L. Estrogen therapy and variable resistance weight training increase bone mineral in surgically menopausal women. J. Bone Min. Res. 6:583-590, 1991. 17. Parfitt A. M., Drezner M. K., Glorieux F. H., Kanis J. A., Malluche H., Meunier P. J., Ott S. M., Recker R. R. (1987). Bone histomorphometry: standardization of nomenclature, symbols and units. J. Bone Min. Res. 2:595-609. 18. Parfitt A. M., Mathews C. H. E., Villanueva A. R., Kleerekoper M., Frame B., Rao D. S. (1983). Relationships between surface, volume and thickness of iliac trabecular bone in aging and in osteoporosis. J Clin Inv 72:1396-1409. 19. Pead M. J., Suswillo R., Skerry, Vedi S., Lanyon L. E. Increased [3H]uridine levels in osteocytes following a single short period of dynamic bone loading in vivo. Calcif. Tissue Int. 43:92-96; 1988. 20. Parfitt A. M. The physiologic and clinical significance of bone histomorphometric data. In: Bone Histomorphometry: Techniques and interpretation. Recker R. R. (ed) (PP.143-223) 1983. 21. Raab D. M., Crenshaw T. D., Kimmel D. B., Smith E. L. A histomorphometric study of cortical bone activity during increased weight-bearing exercise. J. Bone Min. Res. 6:741-749, 1991. 22. Raab-Cullen D. M., Akhter M. P., Dimmel D. B. and Recker R. R. Periosteal bone formation stimulated by externally induced bending strains. J. Bone Min. Res. 9:1143-1152; 1994. 23. Raab-Cullen D. M., Thiede M. A., Petersen D. N., Kimmel D. B., Recker R. R. Mechanical loading stimulates rapid changes in periosteal gene expression. Calcif. Tissue Int. 5:473-478; 1994. 24. Rawlinson S. C. F., Mohan S., Baylink D. J., Lanyon L. E. Exogenous prostacyclin, but not prostaglandin E2, produces similar responses in both G6PD activity and RNA production as mechanical loading and increases IGF-II release in adult cancellous bone in culture. Calcif. Tissue Int. 53:324-329; 1993. 25. Rubin C. T. and Lanyon L. E. Regulation of bone mass by mechanical strain magnitude. Calcif. Tissue Int. 37:411-417, 1985. 26. Smith E. L. and Gilligan C. Dose-response relationship between physical loading and mechanical competence of bone. Bone 18:45S-50S; 1996. 27. Smith E. L., Gilligan C., McAdam M., Ensign C. P. and Smith P. E. Deterring bone loss by exercise intervention in premenopausal and postmenopausal women. Calcif. Tissue Int. 44:312-321, 1989. 28. Tommerup L. J., Raab D. M., Crenshaw T. D. and Smith E. L. Does weight-bearing exercise affect non-weight-bearing bone? J. Bone Min. Res. 8(9): 1053-1058, 1993. 29. Turner C. H., Woltman T. A. and Belongia D. A. Structural changes in rat bone subjected to long-term in vivo mechanical loading. Bone 13:417-422, 1992. 30. Turner C. H. Mechanical loading thresholds for lamellar and woven bone formation. J. Bone Min. Res. 9:87-97, 1994. OUTLINE OF THE PRESENT INVENTION In the instant system, apparatus means is provided for the perfusion and axial mechanical loading of an explanted sample of mammalian trabecular bone which has been prepared in an appropriate manner. During use, a prepared trabecular bone biopsy core is placed within the apparatus and is then loaded mechanically to induce tension and/or compression to the bone matrix. The bone explant perfusion and loading apparatus of the instant system is provided with means for maintaining an environment with stable oxygen, carbon dioxide, nutrients and systemic hormones. Prepared bone biopsies are in the form of trabecular bone cores, 10-12 mm in diameter and 3 to 5 mm thick. These are surgically extracted, under sterile conditions, from suitable long bones of the subject. This procedure is carried out with care and precision, using suitable cutting means and cooling means, to ensure that the resultant bone disk samples are not subjected to temperature rises during cutting and that extreme dimensional accuracy and disk flatness are achieved. Cutting means are in the form of a surgical hand saw used to cut gross samples, a diamond tipped keyhole saw to remove bone cores from the gross samples and an ultra high precision band saw with a diamond tipped bladed, operated in conjunction with jig means, to cut bone disk. Trabecular bone sample disks produced are flat (±100 nm) and have parallel end surfaces (±2-5 cm). Cooling means comprise suitable phosphate buffered saline (PBS) at 6° C. which is used to flood the work piece during cutting. Each trabecular bone sample disk, so prepared, is intended to supply about 3,500,000-11,000,000 cells (based on an estimate of 10,000-20,000 cells per cubic mm of bone (Mundy 1990; Parfitt 1983). Extracted bone disk samples are perfused and maintained with suitable circulating medium, Hepes and fetal calf serum. The apparatus of the instant invention provides means for concurrent mechanical loading of the prepared bone explant sample disk, located within a novel perfusion chamber, in a controlled manner. The maximum compressive strain applied to each sample is 0.5% (5,000 μE) generally at 1 Hz, with the capability of using steeper rise times, if desired. These figures translate to a maximum compression, in each sample, of 20 μm, at a rate of 50,000 μE sec − . Further, the apparatus applies to the bone sample disks, controlled deformations of 200 nm. The apparatus applies forces of up to 800N, at frequencies in the physiological range, of up to 15 Hz and maximum strain rates of between 10,000 μE sec −1 and 50,000 μE sec −1 . These data are appropriate to samples of spongy mammalian bone in which Young's modulus varies between 400 MPa and 1200 Mpa. The apparatus of the instant system also provides an environment in which many factors can be investigated. Because whole tissue is used, bone cells can be studied in a near-natural environment of bone matrix and bone marrow. The apparatus provides means for the user to monitor cellular response but additionally and in a novel manner, to monitor the architecture, strain characteristics and strength of the bone disk and changes therein. The bone explant perfusion and mechanical loading apparatus of the present invention preserves the hard matrix of the bone sample and permits the collection of second messengers and growth factors in the perfusion medium. The instant system thus has many of the advantages of cell culture, whilst retaining the bone matrix encountered in vivo. Means provided within the instant system permit recording of changes in the explanted trabecular bone core sample and further permit the calculation of strain, load and Young's modulus for each such sample. Thus, the instant system permits not only the monitoring of second messengers, cytokines and growth factors but further permits study of how these factors, in conjunction with mechanical loading, will maximize skeletal response to varied stimuli both alone and in combination. In the instant system there are provided perfusion loading apparatus means, power means, control means, computer hardware means, software means and sampling and analysis methods. The perfusion loading apparatus comprises frame means, adjustable biasing pre-loading means, translator loading means, force sensor means and perfusion chamber means. Most components are substantially cylindrical and are accurately machined in corrosion resistant metal, conveniently stainless steel. Frame mounting means are in the form of a relatively massive cylindrical frame housing, comprising a base, a lower frame section, an upper frame section and a cap, each adapted to fit together. These components are secured together with a series, conveniently of 6, partially male-threaded hardened steel bolts which pass through the frame components and are each tightened down with a female-threaded nut. The frame is about 150 mm high and about 80 mm in diameter. The lower part of the frame is substantially solid and has an axial cylindrical hole to accept a ceramic stacked piezo translator which is secured in place by virtue of a close fit in the lower frame and also by screw means through the base. The top part of the frame provides mounting means for adjustable biasing pre-loading means provided by adjustable screw means located axially in and through and the frame cap and secured thereto by threaded means. Within the adjustable biasing pre-loading means there is provided locating and bearing means for force sensor means in the form of an annular quartz crystal force sensor in a precision welded housing. The perfusion chamber assembly is located axially and centrally in the upper section of the frame and comprises a stainless steel bottom bearing cap which provides mounting means for a perfusion chamber body made in durable biologically inert, non-leaching plastics, preferably polycarbonate. A piston, conveniently made in stainless steel, is provided with sealing means in the form of an ‘O’ ring, made from resilient and biologically inert material, preferably neoprene, engages with the upper part of the perfusion chamber body and under the influence of the pre-loading and loading entities, bears down upon a cylindrical explanted trabecular bone sample placed therein. Fluid pathways formed in the perfusion chamber body are disposed so as to ensure that perfusing fluid reaches all parts of the bone sample. Spigots provide connecting means for suitable tube means for delivering perfusing fluid to the assembled perfusion chamber and for collecting effluent from it. The upper and lower components of the perfusion chamber are provided with locating and compression centering means and the assembly is located axially above and upon the translator loading means and directly beneath and in contact with the adjustable pre-load means which drive through push rod and ball bearing coupling means. The piezo translator is provided, via cable connecting means, with a suitable control interface having a microprocessor controlled digital to analogue converter, low voltage driver, controller and power supply, a high voltage amplifier and display unit, all having performance and operating characteristics appropriate to the functional applications of the instant system. The force sensor is provided, via cable connecting means, with a suitable force amplifier having an appropriate power supply and display unit, all having performance and operating characteristics appropriate to the functional applications of the instant system. It will now be apparent that frame means, in co-operation with adjustable biasing pre-load means having force sensor means, translator loading means and perfusion chamber means, as hereinbefore described, constitute perfusion means and instrumented axial press means for the perfusion and mechanical loading of an explanted human or animal bone sample. An explanted trabecular bone sample, prepared as hereinbefore described, is placed within the perfusion chamber, which is then assembled to the frame and loading apparatus. With connections established, power on, and perfusing fluid flowing, the adjustable biasing pre-loading sub-system is adjusted to remove lost motion from and to apply a biasing force to the load train. The biasing force is applied using a large load adjustment knob situated above the frame which drives the adjustable biasing pre-loading means via fine-threaded screw means. A suitable biasing pre-load may also be established using electro-mechanical means via regulator loop means provided in the translator controller. Establishment of a biasing force allows system integrity to be checked. The desired working load or linear translation for the experiment in hand may then be effected using the translator and translator controller. Serial samples of effluent may be collected and assayed for one or more selected factors. Voltage outputs from the translator and charge output from the force sensor are processed and displayed visually. These are used for input to a suitable standard personal computer employing a standard operating system and running a bespoke software program for manipulating data. The program provides software means which produce outputs, via a standard interface, to the system for set-up, configuration, calibration and control of hardware as well as for calculation of relaxation and Young's modulus. Numerical and graphical results may be output to a suitable monitor and printing device connected to the computer. The instant system allows assessment of bone cellular response to specific stimuli, under controlled conditions. An understanding of these mechanisms will allow their manipulation which may possibly lead to the alleviation or control of osteoporosis and other deleterious skeletal changes. The instant system advances the state art in permitting investigators to study physiological responses of bone tissue under specified conditions. The instant system also advances the state of the art in permitting study of human bone biopsies in a controlled environment. It provides means for identifying morphologic changes occurring in different bone diseases and potentially, for the determination of the physiologic and genetic determinants in such diseases. It is thus a first and most important object of the present invention to provide a novel system for continuous perfusion in conjunction with mechanical loading and for collecting and monitoring second messengers, cytokines and growth factors produced by a viable explanted bone sample in order to study skeletal response to varied stimuli both alone and in combination. It is a second important object of the present invention to provide novel means within the instant system for recording changes in thickness of an explanted bone sample during mechanical loading and further for the calculation of strain, load and Young's modulus for each such sample. It is a third important object of the present invention to provide novel apparatus means for concurrent perfusion and axial mechanical loading of an explanted sample of mammalian bone, prepared in an appropriate manner, for an extended period during which the bone is to be kept viable. It is a fourth object of the present invention to provide suitable control and recording means for novel apparatus means for concurrent perfusion and axial mechanical loading of an explanted sample of mammalian bone. The instant system will now be described in more detail in conjunction with the following drawings. DESCRIPTION OF THE DRAWINGS In order that the present invention may be more readily understood, reference will now be made to the following drawings in which: FIG. 1, is a diagrammatic front view of the assembled mechanical and electro-mechanical components of a bone explant perfusion and mechanical loading system, according to the present invention. FIG. 2, is a diagrammatic exploded upper perspective axial view of the mechanical and electro-mechanical components of FIG. 1 . FIG. 3, is a diagrammatic exploded inverted perspective axial view of the mechanical and electro-mechanical components of FIG. 1 . FIG. 4, is a diagrammatic exploded section of the components of the perfusion chamber assembly and a prepared bone sample. FIG. 5, is an underplan view of the perfusion chamber body of the present invention. FIG. 6, is a side section of the assembled components of the perfusion chamber assembly with a prepared trabecular bone sample located therein. FIG. 7, is a schematic diagram of the instant system particularly illustrating electronic control equipment used in conjunction with the electro-mechanical equipment of the present invention. DESCRIPTION OF THE PREFERRED EMBODIMENTS With general reference to FIGS. 1-7, there is described a preferred embodiment of a novel, combined perfusion and mechanical loading system for explanted bone, generally designated by the numeral 10 . Referring to the FIGS. 1-7, there are shown the principal assemblies of a perfusion and mechanical loading system 10 , comprising metal frame housing means 100 , adjustable biasing pre-loading means 400 , translator loading means 600 , force sensor means 700 , perfusion chamber means 800 , and electronic control means 900 . As may best be seen with reference to FIGS. 2 and 3, metal frame housing means 100 , are in the form of a substantially cylindrical and relatively massive metal frame, preferably made from solid stainless steel and comprising several accurately machined parts. In this embodiment, a frame base 102 , is substantially circular, conveniently having a diameter and general thickness of 78 mm and 15 mm, respectively. Additionally, frame base 102 , has a substantially circumferential rim 104 , conveniently 4 mm wide and extending somewhat less than 5 mm above general upper surface 106 , of 102 . A central circular through hole 108 , in frame base 102 , conveniently 4 mm in radius, has a countersunk recess 110 , on the under surface 112 , of 102 . A central circular recess 114 , provided in upper surface 106 , of frame base 102 , is co-axial with central hole 108 . Circular recess 114 , has a radius and depth conveniently of 12.5 mm and 5 mm, respectively. A parallel sided recess 116 , conveniently 9 mm wide, extends from recess 114 , to the edge 118 , of frame base 102 , interrupting rim 104 , thereof. The depth of parallel sided recess 116 , is substantially similar to that of circular recess 114 . A plurality of clearance fixing holes, 120 - 130 , have centers disposed at equal angles around a pitch circle 132 , shown as a center line, conveniently of radius 25 mm and concentric with 108 . Fixing holes, 120 - 130 , are each adapted by the provision of a counter bore 134 - 144 , on under surface 112 , of frame base 102 . Frame base 102 , is provided with a counter bore 146 , disposed in its upper surface 106 , having a diameter and depth, conveniently, of 6 mm and 7 mm, respectively. Counter bore 146 , is centered on a line bisecting the angle between holes 124 ; 126 , substantially midway between pitch circle 132 , and rim 104 of 102 . A lower frame section 148 , is substantially of similar radius to frame base 102 . Lower frame section 148 , conveniently has a height of 61 mm and is provided with substantially circumferential right angled rebates of 5 mm, top and bottom, indicated at 150 and 151 , respectively. A central circular through-hole 152 and a series of fixing holes 154 - 164 , are provided which have substantially similar diameters to and relative spatial dispositions corresponding with 114 and 120 - 130 , respectively, of frame base 102 . Similarly, a parallel sided recess 166 and a counter bore 168 , are provided in the underside 170 , of lower frame section 148 , having dimensions and positions which correspond with 116 and 146 of 102 . Counter bore 168 , is adapted by the provision of female thread means 172 . A locating stud 174 , adapted by the provision of partial male thread means 176 , threadedly engages with the female thread means 174 , of counter bore 168 . Lower rebate 151 and locating stud 174 , of lower frame section 148 , together with rim 104 and counter bore 146 , of frame base 102 , constitute adaptations for mutual location and secure positioning means. An upper frame section 178 , has a height, conveniently, of 65 mm and is of similar radius to 102 and 148 . Additionally and also similarly, upper frame section 178 , is provided with a series of fixing holes 180 - 190 , which have substantially similar diameters to and relative spatial dispositions corresponding to 154 - 164 , of lower frame section 148 and 120 - 130 , of frame base 102 . A central circular hole 192 , has a diameter, conveniently, of 37 mm. Observation means are in the form of a parallel sided fenestration slot 194 , conveniently 30 mm wide, extending from 192 , through wall 196 , of upper frame section 178 , from its upper surface 198 , to a depth, conveniently, of 50 mm. Fenestration slot 194 , is symmetrically disposed about fixing hole 184 and has a deep chamfer 200 , extending out to a width, conveniently, of 53 mm at the outer surface 202 , of 178 . Lower surface 204 , of upper frame section 178 , is adapted by the provision of a machined recess 206 , to form a rim 208 , so dimensioned as to constitute engagement location means for positioning upon upper rebate 150 , of lower frame section 148 . A further parallel sided slot 210 , has substantially the same width and is in a corresponding position to recesses 116 ; 166 , of frame base 102 and lower frame section 148 , respectively. Slot 210 , extends from central hole 192 , through wall 196 , to outer surface 202 and extends to a depth, conveniently of 25 mm, from upper surface 198 , of the 178 . A counter bore 212 , in upper surface 198 , of 178 , is provided with female thread means 214 . A locator stud 216 , is provided with partial male thread means 218 , for threaded engagement with female thread means 214 , of counter bore 212 . Counter bore 212 and locator stud 216 , are of substantially similar dimensions and are similarly spatially disposed to 168 and 174 , of lower frame section 148 . Upper surface 198 , of upper frame section 178 , has a circumferential rebate 220 , substantially similar to 150 , of lower frame section 148 . A frame cap 222 , has a similar diameter and thickness to frame base 102 . Frame cap 222 , is provided with fixing holes 224 - 234 , substantially similar in diameter and disposition to, those of frame base 102 . Fixing holes, 224 - 234 , are each adapted by the provision of a counter bore 236 - 246 , in upper surface 248 , of 222 , substantially similar to 134 - 144 , in 112 , of frame base 102 . Additionally, frame cap 222 , has a rim 250 , disposed circumferentially about its lower surface 252 , substantially similar to rim 104 , of frame base 102 . A counter bore 254 , is adapted and positioned to receive locator stud 216 , of upper frame section 178 . Rim 250 and counter bore 254 , of frame cap 222 , together with rebate 220 and locating stud 216 , of upper frame section 178 , constitute adaptations for mutual location and secure positioning means. A central circular through-hole 256 , adapted by the provision of female thread means 258 , has a diameter, conveniently, of 30 mm. A central circular recess 260 , in upper surface 248 , of frame cap 222 , is conveniently 40 mm in diameter with a filleted circumference to its full depth, conveniently of 4 mm. Rim 250 , is partially interrupted by a shallow cut out 262 , so sized and positioned as to correspond with fenestration slot 194 , of 178 , in order to extend observation means in metal frame housing 100 . Metal frame housing 100 , is assembled by mutually locating and positioning 102 , 148 , 178 and 222 and securing them together with mutual threaded securing means. These means are in the form of a series of bolts, preferable made in hardened steel and indicated in FIGS. 2 and 3, by example bolt 264 . Each bolt has a male-threaded portion, shown on 264 , at 266 and extending through corresponding fixing holes 120 - 130 , in 102 ; 154 - 164 , in 148 ; 180 - 190 , in 178 and 224 - 234 , in 222 . Bolt heads, indicated on 262 , by example bolt head 268 , locate in corresponding counter bores 236 - 246 of 222 . Bolts, exemplified by bolt 264 , are each secured by threaded nut means, in the form of nuts, preferably made of hardened steel and indicated in FIGS. 2 and 3, by example nut 270 , such that each is located in corresponding counter bores 236 - 246 of frame cap 222 . Nuts and bolts are evenly tightened using a torque wrench. Adjustable biasing pre-loading means 400 , are in the form of a plurality of substantially cylindrical, accurately machined components preferably made of stainless steel. A housing 402 , is conveniently 37 mm in maximum diameter and of 65 mm height. A lower portion 404 , is conveniently, 15 mm in height with a diameter of 30 mm and is adapted by the provision of male-threaded means 406 , to threadedly engage from above with female thread means 258 , of axial circular hole 256 , of frame cap 222 . An axial lower counter bore 408 , having a diameter conveniently of 20 mm, extends upwards from lower surface 410 , of 402 , conveniently for 30 mm. A through hole 412 , conveniently of 4 mm diameter, extends radially through the wall 414 , of lower male-threaded portion 404 , of 402 , into lower counter bore 408 , thereof and is adapted by the provision of female-threaded means 416 , throughout its length. An axial upper counter bore 418 , having substantially the same diameter as 408 , extends downwards from upper surface 420 , of 402 , for 15 mm. A further and larger axial upper counter bore 422 , in 420 , conveniently having a diameter of 27 mm and a depth of 10 mm, is adapted by the provision of female thread means 424 . Upper counter bores 418 ; 422 and lower counter bore 408 , are united by an axial circular through hole 426 , conveniently 10 mm in diameter. An adjusting axial screw 428 , conveniently has a length of 60 mm. A lower portion 430 , has a diameter conveniently of 10 mm and a length of 25 mm and is adapted by the provision of male thread means 432 , conveniently extending 22 mm, the remaining 3 mm, indicated at 434 , being plain. An upper plain portion 436 , of 428 , conveniently has a diameter of 10 mm and a length of 30 mm. Upper, plain portion 436 , is provided with radial plain blind hole 438 , conveniently situated half way along its length and having a diameter of 4 mm and extending 2 mm in depth. A central plain portion 440 , of 428 , conveniently has a diameter of 20 mm and a length of 5 mm. A locking collar 442 , is conveniently 27 mm in diameter and 10 mm in depth and is adapted on its outer surface 444 , by the provision of male thread means 446 , to provide threaded engagement means for receival by the female-threaded means 424 , of counter bore 422 , of pre-load adjuster housing 402 . Locking collar 442 , is further adapted by the provision of an axial through hole 448 , having a diameter sufficient to provide sliding engagement means for the upper plain portion 436 , of pre-load adjusting axial screw 428 . Upper surface 450 , is provided with a pair of counter bores 452 ; 454 , disposed on the same diameter of 442 , either side of and equidistant from axis 456 . Counter bores 452 ; 454 , each conveniently have a diameter of 3 mm and extend to a depth of 5 mm and constitute tightening drive means for 442 . A knob 458 , conveniently has a diameter of 60 mm and a maximum depth of 20 mm. An upper portion 460 , having the full diameter of 458 , is conveniently 10 mm deep and has a knurled outer surface 462 . A lower portion 464 , conveniently has a diameter of 30 mm. Knob 458 , is adapted by the provision of an axial through hole 466 , having a diameter such as to provide, in co-operation with upper plain portion 436 , of load adjusting axial screw 428 , easy push-fit means. Knob 458 , is further adapted by the provision of a slot 468 , conveniently 12 mm wide and extending through the full depth of upper portion 460 . Slot 468 , has a semicircular inner margin 470 , adapted by the provision of a small central hole 472 , extending radially into central through hole 466 and is further adapted by the provision of female thread means 474 , for the receival of a grub screw 476 . Grub screw 476 , constitutes engagement locking means between pre-load adjuster knob 458 and pre-load adjusting axial screw 428 . An actuator 478 , has external dimensions such that, in cooperation with lower counter bore 408 , of pre-load adjuster housing 402 , these elements provide fully engageable sliding push-fit means. The upper surface 480 , of 478 , is provided with an axial counter bore 482 , conveniently having a depth of 23 mm, adapted by the provision of internal female thread means 484 , for threaded engagement with lower male-threaded portion 430 , of pre-load adjusting axial screw 428 . Counter bore 482 , is further adapted, at its lower end 486 , by the provision of countersink means, the position of which is indicated by arrow 488 , in FIG. 2 . Countersink means 488 , are for the receival of the head 490 , of countersunk screw means 492 , conveniently of 4 mm diameter. Lower surface 494 , of 478 , is provided with a central counter bore 496 , conveniently having a diameter of 6 mm and which communicates with counter bore 482 . External surface 498 , of 478 , is provided with vertical, parallel sided groove 500 , of such a width and depth as to co-operate, in the assembled condition, with a blunt-nosed grub screw 502 , provided with partial male thread means 504 , received threadedly in and through hole 412 , of housing 402 , to provide engagement guiding and rotation restraining means. It will now be understood that with 502 , engaged in 500 , the latter will be prevented from rotating when knob 458 , is used to turn adjusting axial screw 428 and that instead, actuator 478 , will be driven up or down, within lower counter bore 408 , of housing 402 , according to the direction in which knob 458 , is turned. It will be further understood that this arrangement together with the threaded engagement between internal female thread means 484 , of counter bore 482 , of 478 , with male-threaded portion 430 , of pre-load adjusting axial screw 428 , constitute drive means for setting or altering biasing preload. A push rod 506 , has an upper portion 508 , so sized as to co-operate, slidingly, both with counter bore 496 , in lower surface 494 , of actuator 478 and also with force sensor 700 , hereinafter described. Upper portion 508 , of 506 , is adapted by the provision of a central counter bore 510 , provided with female thread means 512 , for the threaded receival of countersunk screw means 492 , which is introduced down upper counter bore 482 , actuator 478 . Push rod 506 also has a lower portion 514 , having a greater diameter than 508 . Lower surface 516 , of 514 , is provided with a small peripheral chamfer 518 and is adapted by the provision of a central, substantially hemispherical, recess 520 , providing receival means and compression, locating means for a ball bearing 522 , conveniently having a diameter of 6 mm. Ball bearing 522 , provides part of means for transmitting and centering loads applied to perfusion chamber means 800 , hereinafter described. Upper surface 524 , of lower portion 514 , provides shoulder bearing means for push rod 506 , against force sensor 700 . Translator loading means are in the form of a ceramic stacked piezo translator 600 , incorporating multiple strain gauge means and having a maximum translational range of 40 μm. Piezo translator 600 , may conveniently be a commercial product such as a P-239.30, incorporating an optional module P-177.10, having four strain gauges (Physik Instrumente GmbH, Germany). Piezo translator 600 , is substantially cylindrical in form and has a base portion 602 , for which circular recess 114 , in upper surface 106 , of frame base 102 , provides gentle push-fit location means. Translator base portion 602 , has a central counter bore 604 , which is provided with female thread means 606 . Securing means between 600 and 102 , comprise countersunk male-threaded screw means 608 , having a male-threaded shank 610 . Shank 610 , passes through central hole 108 , in frame base 102 and engages, threadedly, with 606 . Countersunk head 612 , of 608 , is tightened against countersunk recess of 110 , of under surface 112 , of 102 . Piezo translator 600 , has a main body portion 614 , having such a diameter that it engages central hole 152 , of lower frame section 148 , with a sliding push-fit. Main body portion 614 , is of such a height that, when fully engaged in assembled metal frame housing 100 , its upper surface 616 , is substantially level with the bottom of fenestration slot 194 , of upper frame section 178 . Upper surface 616 , of piezo translator 600 , is adapted by the provision of an axial counter bore 618 , adapted by the provision of female thread means 620 , for the threaded receival of a small, substantially cylindrical drive pin 622 , having a shank 624 , provided with male thread means 626 . Drive pin 622 , may conveniently be a commercial product such as a P-239.95 (Physik Instrumente GmbH, Germany) described by the manufacturer as a ‘top piece’. Body portion 628 , of drive pin 622 , constitutes boss mounting means for an upper portion 630 , which is substantially hemispherical and has a diameter conveniently the same as ball bearing 522 , of push rod 506 . It is to be understood that 630 , provides the remainder of means, hereinbefore described with reference to ball bearing 522 , means for transmitting and centering loads applied to perfusion chamber means 800 , hereinafter described. When assembled to metal frame housing 100 , connecting means for piezo translator 600 , in the form of cable means are so disposed as to lie in parallel sided recesses 116 , in upper surface 106 , of frame base 102 and 166 , in lower surface 170 , of 148 , providing aperture means constituting access means for cable means to connectors 632 and 634 . Piezo translator 600 , is provided, via connectors 632 and 634 , with electronic control means 900 , all having performance and operating characteristics appropriate to the functional applications of the instant system best seen and described hereinafter, with reference to FIG. 6 . Force sensor means are in the form of a quartz crystal force sensor 700 , housed in an extremely rigid, precision welded, substantially cylindrical housing 702 , having dimensions, conveniently, of outside diameter 14.5 mm, inner diameter 6.5 mm and height 8 mm. Force sensor 700 , may conveniently be a commercial product such as a model 9011A device (Kistler AG, Winterthur, Switzerland). Force sensor 700 , has an axial through hole 704 , for smooth sliding engagement with upper portion 508 , of push rod 506 . In the assembled condition, lower surface 706 , of housing 702 , located on 508 , bears directly upon upper surface 524 , of lower portion 514 , which provides shoulder bearing means for push rod 506 . Upper surface 708 , of 702 , is borne upon by lower surface 494 , of actuator 478 . During use of the instant system, compression of 702 , between 506 and 478 , provides reactive force means for operation of force sensor 700 . Force sensor 700 , is provided with connecting means in the form of cable means which, in the assembled condition, pass through parallel sided slot 210 , of upper frame section 178 . Parallel sided slot 210 , constitutes aperture means in 178 , for access means for cable means to cable connector 710 . Cable connecting means extend from cable connector 710 , to electronic control means 900 , best seen in and hereinafter described with reference to, FIG. 6 . Perfusion chamber means 800 , best seen in FIGS. 4, 5 and 6 , comprise three principal, substantially cylindrical components, a bottom bearing cap 802 and a piston 804 , both machined in suitable grades of stainless steel and a perfusion 806 , preferably machined from a block of suitable biologically inert, non-leaching plastics, preferably polycarbonate and provided with connection means for perfusion fluid. The choice of plastics is very important since many materials leach substances which are toxic or lethal to cells. Minimally different embodiments may be made in which 806 , may be made from suitably biologically inert stainless steels. The preferred embodiment confers the advantage, by virtue of fenestration slot 194 , in upper frame section 178 , of observability of perfusion during use of system 10 . Bottom bearing cap 802 , has a lower portion 808 , conveniently having a radius of 25 mm and a depth of 5 mm. Lower surface 810 , of 808 , is adapted by the provision of an axial, substantially hemispherical, recess 812 , providing receival means and compression, locating and centering means for upper hemispherical portion 622 , of drive pin 618 , of piezo translator 600 . Upper portion 814 , of 802 , conveniently has a diameter of 15 mm and a depth of 5 mm. The upper surfaces 816 ; 818 , of 808 and 814 , respectively, are precision ground to flatness and finished by polishing. Upper portion 814 , of 802 , is provided with male thread means 820 . Piston 804 , conveniently has a diameter of 12 mm and a height of 8 mm. Upper surface 822 , of 804 , is adapted by the provision of an axial, substantially hemispherical, recess 824 , providing receival means and compression, locating and centering means for ball bearing 522 , of push rod 506 . Lower surface 826 , of piston 804 , is precision ground to flatness and finished by polishing. Circumferential wall 828 , of 804 , is adapted by the provision of an upper and a lower annular groove, indicated at 830 and 832 , respectively and mutually disposed apart in a parallel manner to upper and lower surfaces 822 and 826 , respectively. Circumferential wall 828 , is finished by micro-fine machining and polishing. Lower annular groove 832 , constitutes an adaptation for the receival of a sealing means in the form of an ‘O’ ring 834 , made from inert resilient sealing material, preferably neoprene. Perfusion chamber body 806 , conveniently has an outer diameter of 25 mm and a height of 15 mm. Lower surface 836 , of 806 , is provided with an axial counter bore 838 , of such a depth and diameter as to provide, in conjunction with suitable female thread means 840 , receival and sealing means for upper portion 814 , of bottom bearing cap 802 . Lower surface 836 , of 806 , is precision ground to flatness and is adapted to cooperate with upper surface 816 , of lower portion 808 , of bottom bearing cap 802 and suitable biologically inert non-leaching adhesive means, to provide additional sealing means between the two components. Lower surface 836 , of 806 , is further adapted by the provision of a machined annular channel 842 , conveniently having a semi-circular cross-section of 2.5 mm diameter and lying on a pitch circle, conveniently 19 mm in diameter and indicated with center line at 844 , in FIG. 5 . Two small counter bores 846 and 848 , having the same diameter as 842 , are each centered on the intersection of a diameter of 842 and center line 844 , one either side of central axis 850 . Small counter bore 846 , extends to a depth somewhat less than that of the main lower axial counter bore 838 . A radial hole 852 , having the same diameter as counter bore 846 , extends through wall 854 , of 806 , so as to meet 846 , at right angles, forming substantially continuous lumen means. An upper axial counter bore 856 , is so sized and adapted that it may receive an explanted trabecular bone sample 858 , prepared as hereinafter described, as a sliding fit and also may engage the greater part of piston 804 , as an easy push fit. Lower annular groove 832 , of piston 804 and ‘O’ ring 834 , located therein, are particularly included in the engagement between piston 804 and perfusion chamber body 806 . Inner surface 860 , of upper axial counter bore 856 , is adapted by the provision of a parallel sided, annular channel, 862 , conveniently 4 mm wide and about 2 mm deep. The position of 862 , is such that it substantially surrounds the outer margin or wall 864 , of explanted trabecular bone sample 858 , when this is inserted in 856 , of 806 . Second small counter bore 850 , extends upwards into 806 , to a depth somewhat greater than the depth of main lower axial counter bore 838 , such that it terminates at a point substantially level with the mid point of the height of inserted explanted bone sample 858 . A diameter hole 866 , having the same diameter as small counter bores 846 ; 848 , extends through wall 854 , of 806 , following the line of radial hole 852 , hereinbefore described, intersecting annular channel 862 and also intersecting second small counter bore 848 , at right angles, at the limit of its depth, forming further substantially continuous lumen means. Spigots 868 and 870 , conveniently fabricated in stainless steel, are adapted to engage, respectively, with radial hole 852 and that portion 872 , of diameter hole 866 , which lies on the same side of 806 , as 852 , with a forced, sealing, press-fit. A small cylindrical plug 874 , of the same material as 806 , is adapted to engage with diameter hole 866 , on the opposite side to radial hole 852 , with a press-fit in conjunction with suitable biologically inert, non-leaching, adhesive means to provide sealing means between the two components. Plug 874 , is of such a length that it extends up to but does not substantially encroach into, second small counter bore 848 . Annular channel 842 , of lower surface 836 , of perfusion chamber body 806 , in cooperation with the upper surface 816 , of lower portion 808 , of bottom bearing cap 802 and adhesive sealing means; small counter bores 846 and 848 , annular groove 862 , of upper counter bore 856 , radial hole 852 and plugged portion 876 , of diameter hole 866 , constitute substantially continuous fluid pathway means for perfusing fluid. Spigots 868 and 870 , constitute connecting means for suitable tube means in the form of tubes conveniently made of silicone rubber and indicated at 878 and 880 , for delivering perfusing fluid to the assembled perfusion chamber and for collecting effluent from it for monitoring and analysis. It will now be understood that the substantially cylindrical elements of frame means, in co-operation with adjustable biasing pre-loading means having force sensor means, translator loading means having electronic control means and connection means and perfusion chamber means having connecting means for perfusing fluid, as hereinbefore described, constitute instrumented perfusion and axial press means for the perfusion and mechanical loading of an explanted trabecular bone sample. Performance and Function of Piezo Translator, Force Sensor and Electronic Control Means With particular reference to FIG. 6, as well as continuing reference to FIGS. 2 and 3, electronics control means 900 , comprises a rack 902 , in which are mounted several major components. A 220V AC power supply 904 , also houses a display module 906 , which gives readings of high voltage or compression values. A high voltage amplifier 908 , provides the high operating voltage (−1000V) to drive piezo translator 600 . A controller module 910 , includes a compression signal amplifier (not seen) and regulator loop (not seen), to force piezo translator 600 , to a required position, within its translational range of 40 μm, corresponding to a given value of high voltage or compression. This range is satisfactory for applications involving explanted bone samples in the instant invention. Controller module 910 , may conveniently be a commercial product such as E-255 PZT Interface and Controller (Physik Instrumente GmbH, Germany) which incorporates a digital-analogue converter (DAC). Controller module 910 , is linked by cable means (not seen) to low voltage driver and controller 912 , which may conveniently be a commercial product such as LVPZ Driver and Controller E-809 (Physik Instrumente GmbH, Germany). Controller module 910 , is also linked by cable means (not seen) to a force signal amplifier 914 , which is a charge amplifier for amplifying output from force sensor 700 . A personal computer 916 , is equipped with a microprocessor of at least 386 rating and is provided, internally, with an additional plug-in card (not seen) which provides a control interface between an analogue-digital converter (ADC) and DAC of 910 . Cable 918 , connects the additional plug in card of 916 , to a compression signal amplifier output provided on low voltage driver and controller 912 . Cable 920 , connects the additional plug in card of 916 , to force signal amplifier 914 . Cable 922 , connects between a communications port COM-1 (not seen) of 916 , to the digital-analogue converter of 910 . Cable 924 , connects between a communications port COM-2, (not seen) of personal computer 916 and a mouse 926 . Personal computer 916 , is also provided with a local printer terminal port (not shown) for the connection of a suitable printer (not shown). Personal computer 916 , is also equipped with a graphics monitor 928 , functioning to EGA, VGA or higher standard to which it is connected by a monitor cable 930 . A suitable operating system, such as DOS™ 3.2 or higher or Windows 3.1™ or Windows 95™, is installed on personal computer 916 , together with custom software which provides means for coordinating and calibrating the electro-mechanical elements of the system as well as for collecting, collating and displaying data and making calculations thereon and displaying the results thereof. In FIG. 6, frame means 100 and adjustable biasing preloading means 400 , are shown in side view to reveal connecting means for cable means. Cables 932 and 934 , connect high voltage amplifier 908 , to piezo translator 600 , at connectors 632 and 634 , respectively. Cable 936 , connects force sensor 700 , to controller module 910 , at connector 710 . Piezo translator 600 , incorporates four strain gauges (not seen) attached internally to the ceramic stack and arranged in a full Wheatstone bridge circuit. The multiple strain gauge arrangement may conveniently be in the form of an optional commercial module P-177.10 (Physik Instrumente GmbH, Germany). In conjunction with controller module 910 , the bridge arrangement allows a positioning accuracy of 0.2% of the nominal expansion of piezo translator 600 , to be achieved. Force sensor 700 , is a quartz crystal force sensor for measuring dynamic and quasi-static forces, having a range of 15 kN, a very high resolution of 0.01N under any pre-load, sensitivity of ≈ — −4.3pC/N, modulus of 3.6 GPa and very high rigidity of ≈ — −1.8 kN/μm. These characteristics are satisfactory for applications involving explanted bone samples in the instant invention. Experiment 1—Calibration And Validation of Loading Elements Of The System The instant system was validated and characterized by the following methods: a. determination of any errors in the system b. identifying deformation accuracy, force application, frequency of loading and calculation of E (Young's modulus) on known materials and determining the physical compliance in the system. Calibration and validation was accomplished by comparing nondestructive test results of the instant mechanical loading, translator and force sensor elements of the instant system to identical tests run on an MTS (Bionix) servohydraulic test machine. Homogeneous materials, with moduli that span the expected range of cancellous bone, were used (e.g. nylon, aluminum, teflon). These materials had strain gauges applied to the vertical surfaces. Strain was monitored on the same materials in both systems and the results were compared. In addition, a precision extensometer was placed between the platens on the MTS machine to provide specimen deformation, as well as load and thus compute the strain. The current required to achieve similar deformations, strains and loads was recorded. The systems were compared with ramp and sinusoidal wave forms. Hysteresis was noted together with time dependent responses in the materials and test system. The system was the materials and test system. The system was tested quasistatically and at increasing frequencies up to 10 Hz (a functional limit for the MTS system). The system was also tested throughout the range of functional deformation rates available with the piezo crystal translator. Similar specimens were taken to failure and the total material behavior curves of the MTS system and the instant system using the piezo crystal translator, were compared. Correlation of a very high order was established, validating the prospective deployment of the novel mechanical loading system, in conjunction with the novel perfusion means of the instant invention in explanted trabecular bone samples. It was determined that the mechanical and electro-mechanical elements of the instant system are capable of applying controlled deformations, accurate to 200 nm, and applying forces of up to 800 N, at frequencies in the physiological range of up to 15 Hz and maximum strain rates of between 10,000 μE sec−1 and 50,000 μE sec−1. Young's modulus for trabecular bone varies from E=400 MPa to, typically, E=1200 MPa in the adult pig. Experiment 2—Perfusion Preparation of Explanted Bone Samples Features considered when determining the optimal size of the bone sample for the instant system were: 1. The practicality of using cow and pig bone samples in the first instance and the feasibility of using human bone samples, subsequently, having the same dimensions. 2. The volume of bone and means for achieving adequate perfusion through it. 3. The amount of tissue which would be necessary to produce the desired biochemical markers, in quantities sufficient to make the required measurements. The selection of pig and cow trabecular bone was based on earlier studies by the present authors and other workers. In particular the studies of an associate, Dr. Kit Mui Chiu whose observations were recorded in a doctoral thesis at the University of Wisconsin, presented in 1996 and entitled “The effect of carnitin dehydroepiandrosterone sulfate on young senescent osteoblast-like cells”, were important. In these studies pig osteoblasts were kept viable, in culture, for 68 days. Careful consideration of these findings and other prior art, led to the conclusion that, in a suitable novel system, providing continuous perfusion means and suitable loading means, viability might be maintained for a worthwhile period of study which could be up to 14 days or more. The bone cores for our experiments were obtained from the trabecular bone of distal ulnae or femurs of 2 to 3 year old cows or femora or humeri of 2 to 3 year old pigs. Under sterile conditions throughout, the limb is first excised and then a 2.5 cm×2.5 cm×4.5 cm (proximo-distal dimension) sample of trabecular bone is cut from the central region of the proximal or distal metaphysis of the bone with a surgical hand saw and the proximal end is marked. The specimen is visually inspected under a dissecting microscope at 10× to assure that no growth plate scars are present. Following isolation of the gross sample, 6×5 mm thick subspecimens are cut from it, under running sterile PBS at room temperature, using a band saw having a diamond tipped blade (Exact, Germany). Six bone core disks are then drilled in the proximo-distal direction, under sterile PBS, from each of the sub-specimens, using a 10 mm or 12 mm diamond tipped keyhole drill (Exact, Germany). The 6 bone cores from each 5 mm sub-specimen are randomized. Each bone core disk is immediately marked on the proximal surface and placed in serum free medium for 20 minutes prior to placing it in the perfusion chamber apparatus of the present invention. Each sample is placed in the perfusion chamber such that it will be loaded from proximal to distal. The sample is then allowed 48 hours in the perfusion chamber in order to adapt, prior to any intervention. Thus, all experiments conducted using this protocol extend over 16 days, comprising 2 days for core adaptation and 14 days of intervention. Using this method, bone disks may be cut with the necessary extreme precision to a flatness of ±0.2 microns and a parallelism of ±0.1 microns. The dimensions were selected in order to produce samples of a practical size for perfusion and in order to supply between 3,500,000-11,000,000 cells, based on an estimate of 10,000-20,000 cells per millimeter cube of bone (Mundy 1990; Parfitt 1983), which was considered sufficient to provide an adequate yield of markers for study. Disk samples of trabecular bone, prepared according to the method immediately hereinbefore described, were perfused and maintained with circulating medium. The medium used was Ham's F10 containing 1%-5% FCS, 2 mg glutamine, streptomycin and penicillin G at 50,000 U/1, vitamin C 10 mg/ml, 0.12 g/l of NaHCO 3 and 10 mM Hepes. The medium was maintained at 37° C. and a pH of 7.1-7.3 for the total 14 days of the perfusion. The perfusion rate was 0.1 ml/minute and the medium was perfused using a 12 channel pump (Ismatec). The medium was changed at 12 hour intervals. The pH, PCO 2 and PO 2 were measured hourly for the first 5 hours then 12 hourly thereafter. A series of FCS batches was tested for biological effect on the trabecular bone cores using alkaline phosphatase, cell viability and osteocalcin production. A sufficient quantity was retained from the most suitable batch of FCS to maintain a reproducible medium for the performance of the experimental program contemplated by the investigators. It is important to note that frozen FCS (−80° C.) has a maximum storage life of 3 years. The flow rate through the explanted bone sample must be fast enough to maintain cell viability but not so fast that a shear force greater than 3 dynes/cm 2 is induced. When the flow rate is too slow, cells are inadequately oxygenated and lactate builds up. When the flow rate is too fast, the shear force, itself, causes increases in PGE2 and IGF-1. The flow rate of 0.1 ml/minute selected was determined as optimal by prior experiment with differing flow rates in order to provide sufficient effluent medium volume for sampling and analysis of PGE2, cAMP and IGF-1 and also in order to maintain PO 2 and PCO 2 . PO 2 was monitored at each flow rate in these experiments to ensure adequate oxygenation of the cells in the bone explant perfusion/loading system. Experiment 3—Injury Response Time (Establishment of Rest Period) Cells placed in culture require time to adapt to their changed environment and this time period varied with the type of cell and the type of research we conducted. The necessary rest period for explanted trabecular bone samples in the instant system was determined. In our preliminary experiments, trabecular bone samples were perfused with culture medium plus 10% FCS. Under these conditions, IGF-1 increased from 5 to 14 hours and appeared to decline in the 15th hour, at the time the experiment was terminated. Based on those preliminary data, a rest period of at least 48 hours was accepted as appropriate for IGF-1 to return to baseline level, before any intervention (mechanical loading, hormones, etc) was imposed on the bone explant organ culture. However, the adaptation time required was then documented over a series of full 24 hour periods to determine when the cells had recovered from the surgical trauma in order to determine the stable baseline condition from which intervention could be started. Studies have provided the equivalent data for each of the second messengers, IGF-1 and certain other growth factors. Experiments to Investigate Cell Viability And Biomarkers under Varying Conditions Experiments were designed to investigate a variety of load magnitudes and frequencies, growth factors and applied active substances. Specifically it was considered necessary to provide for the investigation of markers including the release of prostaglandin E2 (PGE2), cyclic-AMP (cAMP), inositol 1,4,5-trisphosphate (IP3) and insulin-like growth factor (IGF-1), in the perfusion effluent from explanted bone samples. These entities were to be studied during responses to stimuli including varying conditions of mechanical load and further, under the influence of biochemical stimulus with hormones, growth factors or drug substances. The markers, produced by stimuli, immediately hereinbefore described, are important in the regulation of bone modeling and remodeling, at every age and nutritional level, in the adaptive response of the skeleton to such challenges. Cell Viability a. Cell viability in samples of cow trabecular bone was determined at rest, at a maintenance load (the load at which the bone neither atrophies or hypertrophies) and at microstrains which ranged from 500 to 5000. b. Oxygen utilization of the bone explant perfusion/loading model was determined at rest, at varied flow rates, at a maintenance load and at microstrains ranging from 500 to 5000. Having established the flow rate limits for the instant perfusion chamber, experiments were conducted to verify cell viability. Percent viability at various time intervals was assessed in order to determine the number of cells still alive at any given time. Two methods are commonly used to assess cell viability in cell culture. Alamar Blue Assay indicates succinate dehydrogenase activity in the cells. It incorporates an oxidation-reduction indicator that causes the Redox indicator to change from oxidized (non-fluorescent, blue) form to reduced (fluorescent, red) form in response to the cell metabolism in the culture medium. This assay is a general indicator of the metabolic function of the system but it does not allow quantification of cell viability, that is, calculating the percentage and distribution of viable cells. The use of MTT (sigma, 3-[4,5-Dimethylthiazol-2-yl]-2,5-diphenyltetrazolium bromide), does permit measurement of the number or percentage of live cells. In this assay, active mitochondrial dehydrogenases convert the water soluble MTT into an insoluble purple formazan by cleaving the tetrazolium ring. Cells with intact mitochondria will show a dark reddish/purple stain when tissue is viewed under a light microscope. Thus, mitochondrial staining is indicative of live cell function at the time MTT is administered. Since MTT is toxic to the cells, it can be used only at the end of an experiment. We used this method to determine the viability of cells after 14 days. At the end of perfusion runs, samples were perfused with MTT (30 mg/ml) for 6 hours then stabilized to 40° and sectioned to between 100 and 180 microns thick, using a diamond saw (EXACT, Germany), in order that viability throughout the sample could be investigated. Somewhat less than 5% of the cells in the cow trabecular bone core, taken from the distal ulnae of 24 month old cows, died because of the surgical extraction and disk preparation procedure and of the remaining cells, more than 95% remained active throughout the 14 day studies. Histological Assessment Before loading the explanted bone samples, it was necessary to verify whole explant tissue viability over time. MTT (30 mg/ml) was used as a cell viability marker. In one series of our cell viability experiments, four 14-day runs were conducted. In the first two runs, sample cores were processed with MTT every two days after a baseline core had been run for 8 hours and then removed. All cores were compared to the baseline core. 12 bone cores were perfused at a rate of 0.1 ml/minute. A baseline positive control viability sample was obtained 8 hours after the start of the experiment by perfusing a core with 30 mg/ml of MTT for 6 hours. The baseline sample and all other samples were perfused with serum free medium for the first 24 hours in order to collect 1 ml of medium for IGF-1 and PGE2 analysis. At the end of the 6 hour perfusion with MTT, the bone core was maintained at 4° C. at which time sections were cut using a diamond band saw. Sample sections were cut to a thickness between 100 and 180 microns in order that cell viability could be determined throughout the sample. The base line sample (8 hours) was used as a positive control for viability. The number of viable cells in the 14 day sample showed no difference when compared to the positive control which had 95% viable cells. The sample sections taken from the top to the bottom of the sample demonstrated no difference in the number of cells showing the presence of MTT and the centers of all of the cores were found to be fully stained. However, there were a few trabecular areas that demonstrated cell death with no MTT present. It was felt that the diamond tipped keyhole drill used to excise the bone samples may have resulted in some damage in the outer few trabecular segments, resulting in tissue damage and cell death. It is clear from our results in this study that the bone cores obtained using this method and using the perfusion chamber apparatus of the present invention, can be maintained in a viable for 14 days. Bioassay Medium from the perfusion chamber to be used for the bioassay was sampled at varied time intervals according to the biomarker we chose to investigate. Pig osteoblasts obtained from Crenshaw (U of WI Madison) were characterized by alkaline phosphatase, collagen type 1 and the ability to produce bone nodules. Cells were plated out in 96-well, Nunclon, cell-culture grade, assay plates at a density of 45,000 cells per cm 2 in 100 ml per well of one of the following media: Dulbecco's MEM Dulbecco's BGJ (as used for the organ culture) Ham's F-10 HI growth enhancement medium (Gibco) The specific medium was chosen through trial and error depending on the best response of the markers we investigated (e.g. good for alkaline phosphatase and collagen). To the selected basic medium was added 10% FCS, ascorbic acid-2-phosphate at 5 mg/1 plus L-glutamine (or the stable analogue) for the first 24 hours. For the assay, the FCS is reduced from 10% to 1%, for 24 hours before the medium is replaced with medium from the perfusion culture. The control is unused medium used for the perfusion culture. Eight replicate wells were used for each sampling point. The cells were grown for 48 hours and then assayed for growth using the MTT method to measure succinate dehydrogenase activity. The MTT methods were calibrated against a known number of cells in a similar growth state; this was a control experiment using an agar plate and counting the cells with a cell counter. The presence of growth factors released from the perfusion culture were then assayed. Loading In loading experiments, the maximum compressive strain applied was 0.5% (5,000 μE) at 1 Hz sine wave. This equates to 20 cm compression at up to 50,000 μE sec −1 . The bone explant perfusion/loading system we have developed has allowed us to assess bone cellular response to specific stimuli under controlled conditions. An understanding of these mechanisms allows for their manipulation and in turn may lead to the possible alleviation or control of osteoporosis and other skeletal changes which result in the loss of skeletal integrity and function. The instant system provides investigators, for the first time, with effective means to study morphological changes in the skeletal tissue. In addition, the instant system permits the study of the physiological responses of the bone tissue under clearly defined and specified experimental conditions that can be set up to reflect the human activities of daily living and life style. The present invention also for the first time, permits the study of human bone biopsies in a controlled environment. This will not only enable investigators to identify morphologic changes that occur with different bone disease but will also permit the determination of the physiologic and possibly genetic determinants in such conditions. It will be apparent to those skilled in the art that numerous modifications or changes may be made without departing from the spirit or the scope of either the present invention or its method of use. Thus the invention is only limited by the following claims.

A mechanical testing device has a rigid frame and a piezo translator connected to the frame. A Wheatstone bridge is connected to the translator to produce an electrical signal related to the compression of the translator, wherein a sample positioned between the piezo translator and the frame is subjected to loads by the movement of the translator. A sensor detects the force applied to the sample by the piezo translator, and produces a signal indicative of the force. A computer receives the Wheatstone bridge electrical signal and the signal indicative of the force applied to the sample. The computer controls the advancement of the translator to allow the application of precise amounts of compression to the sample.

FIELD OF ART [0001] The field of art relates generally to Carbon nanotubes (CNTs); more specifically, techniques for sorting CNTs of different electronic types. BACKGROUND [0002] A Carbon nanotube (CNT) can be viewed as a sheet of Carbon that has been rolled into the shape of a tube. CNTs having certain properties (e.g., a “metallic” CNT having electronic properties akin to a metal) may be appropriate for certain applications while CNTs having certain other properties (e.g., a “semiconducting” CNT having electronic properties akin to a semiconductor) may be appropriate for certain other applications. CNT properties tend to be a function of the CNT's “chirality” and diameter. The chirality of a CNT characterizes its arrangement of carbon atoms (e.g., arm chair, zigzag, helical/chiral). The diameter of a CNT is the span across a cross section of the tube. [0003] Because the properties of a CNT can be a function of the CNT's chirality and diameter, the suitably of a particular CNT for a particular application is apt to depend on the chirality and diameter of the CNT. Unfortunately, current CNT manufacturing processes are only capable of manufacturing batches of CNTs whose tube diameters, lengths and chiralities are widely varied. The problem therefore arises of not being able to collect CNTs (e.g., for a particular application) of only a particular size range (e.g., length and/or diameter range(s)) and/or electronic property (e.g., metallic or semiconducting) from a batch of manufactured CNTs having widely varied sizes of both metallic and semiconducting CNTs. [0004] CNTs are also known to have poor solubility. Here, owing to van der Waals forces (it is believed), individual CNTs tend to “bundle together” into groups. Thus, when a batch of manufactured CNTs are made to flow in a fluidic stream (such as an aqueous solution), bundled groups of CNTs are observed drifting/flowing through the liquid together. [0005] Success at improving the solubility of CNTs has been reported. For example, Zheng et al. (“DNA-Assisted Dispersion and Separation of Carbon Nanotubes”, Nature Materials 2, pgs. 338-342, 2003) has published a process by which single stranded DNA (ss-DNA) is used to “break-down” a CNT bundle into individual CNTs wrapped in a helical structure of DNA. Here, a CNT that has bonded in some fashion with DNA so as to form a combined structure of DNA and the CNT is referred to as “DNA/CNT hybrid structure”. A DNA/CNT hybrid structure of DNA and a metallic CNT may be referred to as “DNA/metallic CNT hybrid structure”. A DNA/CNT hybrid structure of DNA and a semi-conducting CNT may be referred to as “DNA/semi-conducting CNT hybrid structure”. A “CNT hybrid structure” can be viewed as a CNT that is attached to another substance. [0006] According to the technique taught by Zheng et al., an aqueous solution containing bundles of CNTs is subjected to the presence of ss-DNA. Because the binding energy associated with the coupling of ss-DNA to a CNT is comparable to the binding energy associated with the coupling of CNTs to one other, the application of sonic energy to the solution can create dynamic situations in which an individual CNT that is bundled with one or more other CNTs will reach a lower energy state if the CNT binds with ss-DNA molecules instead of the CNTs associated with its bundle. Because physical systems tend to fall to lower energy states, this prompts the formation of an individual (i.e., non bundled) CNT helically wrapped in ss-DNA. That is, the CNT essentially leaves its bundle in favor of being helically wrapped by ss-DNA. [0007] According to follow-up work reported by Zheng et al. in “Structure-Based Carbon Nanotube Sorting by Sequence Dependent DNA Assembly”, Science 28 Nov. 2003; 302: 1545-1548, a particular sequence of ss-DNA can be made to self assemble into a helical structure that wraps around the surface of an individual CNT. Individual CNTs wrapped by ss-DNA can then be sorted according to their electrical characteristics through anion exchange chromatography. In this manner, individual CNTs having specific “sought-for” electrical characteristics can be collected. FIGURES [0008] The present invention is illustrated by way of example and not limitation in the figures of the accompanying drawings, in which like references indicate similar elements and in which: [0009] FIG. 1 a shows a flow of CNTs being injected into an array; [0010] FIG. 1 b shows the flow of FIG. 1 a emerging on the other side of the array; [0011] FIG. 2 shows an emergence profile for various types of CNTs that flow through an array; [0012] FIGS. 3 a through 3 d show various types of arrays; [0013] FIGS. 4 a and 4 b shows a system for sorting CNTs that contains an array. DETAILED DESCRIPTION [0014] A technique for sorting a batch of metallic and semiconducting CNTs of differing lengths and diameters is to flow them through an arrangement of fixed structures such as an array of posts or studs. Owing to the dynamics by which the CNTs flow through the array as described in more detail below, CNTs of certain size and electronic property type will emerge from the array at a certain time and/or location. FIGS. 1 a and 1 b show a high level perspective. [0015] According to FIG. 1 a , a flow 101 of metallic and semiconducting CNTs of various lengths and diameters are observed flowing into the middle of an “entrance” edge of an array 102 . FIG. 1 b shows the flow of CNTs at a moment of time after that depicted in FIG. 1 a . Importantly, as observed in FIG. 1 b , the smaller CNTs are observed emerging 103 _ 1 , 103 _ 3 from an “emission” edge of the array 102 farther out from the middle of the array 102 than a mixture of larger and smaller CNTs 103 _ 2 . Therefore the act of flowing various CNTs through an array can be used as a basis for separating the CNTs into categories of likeness (e.g., by gathering smaller CNTs at the outer edges of the “emission” edge 102 ). [0016] Two basic separation principles that may simultaneously affect the flow of CNTs through the array include: 1) smaller CNTs (principally based on CNT length rather than CNT diameter) will diffuse farther out along the edge of the array than larger CNTs (as alluded to above with respect to FIGS. 1 a and 1 b ); and, 2) if the array's posts are designed to promote some type of electric field based interaction with the DNA/CNT hybrid structures, semiconducting hybrid structures are more apt to be attracted to the array's posts than metallic hybrid structures resulting in hybrid structures with metallic CNTs emerging from the array farther out along the array edge and earlier in time than hybrid structures with semiconducting CNTs of comparable length and diameter. [0017] The first principle above stems from basic principles of diffusion. Here, larger CNTs (principally, longer CNTs) will diffuse less in comparison to smaller CNTs and upon collision with the posts they will be confined to the vicinity of the center of the array. Here, both smaller and larger CNTs will collide with the array, but smaller CNTs with collide with posts further from the center of the input stream as a result of their larger diffusion; which, in turn, effectively corresponds to the outward surge of the larger CNTs, as a whole, being suppressed by the array posts to a greater extent than the collective outward surge of the smaller CNTs. If the rate of the total flow of CNTs through the array is less forcibly induced (e.g., with a “slow” fluidic flow through the array), it is also possible that smaller CNTs will diffuse out of the array earlier in time than the larger CNTs. [0018] The second principle described above is a consequence of the fact that semiconducting CNTs do not promote the formation of “image” charges on the CNT surface as prevalently as metallic CNTs do. Here, as is well known in the art, metals form image surface charges of a first polarity in the presence of charges of opposite polarity. For example, the presence of a positive charge proximate to a metal will cause negative image charges to appear at the metal's surface. The ability to form such image charges is directly related to the very high mobility of electrons within metals. Because semiconducting materials do not possess electron mobilities as high as metals, semiconducting materials do not form image charges as strongly as a metallic materials. [0019] As such, metallic CNTs are apt to have a greater propensity for producing image charges on their respective surfaces than semi-conducting CNTs. In this regard, an additional sorting mechanism may occur if the CNTs are individually wrapped in DNA as described in the background section of the present application. Specifically, because strands of DNA are negatively charged, DNA/metallic CNT hybrid structures are apt to be neutrally charged (because the underlying metallic CNT will have induced positive image charges that effectively cancel the DNA's negative charge), while, DNA/semiconducting CNT hybrid structures are apt to be negatively charged (because the underlying semiconducting CNT will not have induced sufficient positive image charges to effectively cancel the DNA's negative charge). [0020] As a consequence, if the array's posts are given an electrical charge, DNA/metallic CNT hybrid structures are apt to be unaffected (i.e., there will be little if any attraction/repulsion between the hybrid structure and the array post), while, DNA/semiconducting CNT hybrid structures are apt to exhibit some kind of attraction to or repulsion from an array post depending on the array post's charge (specifically, attraction if the post is charged positively and repulsion if the post is charged negatively). In the case of attraction, the DNA/semiconducting CNT hybrid structures will: 1) progress slower through the array than the DNA/metallic CNT hybrid structures; and, 2) will have their outward surge thwarted to a greater degree than the DNA/metallic CNT hybrid structures. As a consequence, hybrid structures with semiconducting CNTs will emerge from the array later in time and closer to the entrance position of the input flow than hybrid structures with metallic CNTs. [0021] In the case of repulsion (i.e., where the array posts are negatively charged), the opposite separation dynamics are apt to occur. Specifically, the DNA/semiconducting CNT hybrid structures will: 1) progress faster through the array than the DNA/metallic CNT hybrid structures; and, 2) will have their outward surge promoted to a greater degree than the DNA/metallic CNT hybrid structures. As a consequence, hybrid structures with semiconducting CNTs will emerge from the array sooner in time and farther from the entrance position of the input flow than hybrid structures with metallic CNTs. [0022] Moreover, as reported in the work by Zheng, because DNA wrapping geometries change as a function of CNT diameter, the larger the diameter of a CNT within a hybrid structure, the more the CNT will behave like an electric dipole. A CNT in the form of an electric dipole can be viewed as, owing to the presence of an electric field within the CNT, one tip of the CNT being positively charged and the other tip of the CNT being negatively charged. [0023] Here, if the array's posts are positively charged, an electric dipole CNT will position itself so that its negative tip is oriented toward an array post. This effectively corresponds to a form of attraction between the dipole-like CNT and the array post, which, similar to the attractive force described just above, slows the rate at which the CNT can flow through the array and shorten the distance the CNT will emerge outward along the array edge. Because semiconducting CNTs are more apt to exhibit dipole-like behavior than metallic CNTs (because semiconducting CNTs are more dielectric-like than metallic CNTs and are therefore more capable of sustaining an internal electric field), similar separation dynamics may occur in which hybrid structures with semiconducting CNTs emerge later in time and closer to the array entrance position than hybrid structures with metallic CNTs (assuming the CNTs for both kinds of hybrid structures have similar lengths and diameters). [0024] At a high level, the principles described above can be summarized as demonstrating that: 1) principles of diffusion will promote smaller Hybrid structure to emerge farther out along the array emission edge (and perhaps sooner than) larger Hybrid structure; and, 2) if the arrays are “treated” so as to promote some type of attraction or repulsion interaction with specific types of Hybrid structures, the Hybrid structures that exhibit the strongest attraction to the array posts are apt to emerge later in time and closer to the array entrance point than other Hybrid structures, or, contra wise, if the array posts are treated so as to promote repulsion interaction with specific types of Hybrid structures, the Hybrid structures that exhibit the strongest repulsion to the array posts are apt to emerge sooner in time and farther out along the array's emission edge than other Hybrid structures. [0025] FIG. 2 further elaborates on these principles. Specifically, FIG. 2 plots array emergence position as a function of time for a flow of DNA/CNT hybrid structures through an array of positively charged posts where a slow flow rate exists so as to enhance the impact of pure diffusion on the sorting dynamics. Here, consistent with the principles discussed above, specific groups of hybrid structures 201 through 206 are observed as having specific, respective ranges of: 1) time of emission from the array (horizontal axis); and, 2) location along the emitting edge of the array (vertical axis where a position that is “higher up” on the vertical scale corresponds to a position that is farther out along the emitting edge of the array than the position on the entrance edge of the array where the flow of CNTs are introduced). [0026] According to FIG. 2 , longer, larger diameter, semiconducting CNT based hybrid structures 206 generally emerge later in time and closer to the entrance position of the array than the other hybrid structures; and, shorter, smaller diameter, metallic CNT based hybrid structures 201 emerge sooner in time and farther out along the array's emission edge than the other hybrid structures. Here, shorter, smaller diameter metallic CNTs could easily be collected by collecting hybrid structures at time t 1 and position x 2 . Likewise, longer, larger diameter semiconducting CNTs could easily be collected by collecting hybrid structures at time t 2 and position x 1 . [0027] Notably, metallic CNT based hybrid structures 201 through 203 tend to emerge sooner than semiconducting CNT based hybrid structures 204 , 205 , 206 of comparable length and diameter. Amongst the metallic CNT based hybrid structures 201 , 202 , 203 , the hybrid structures with shorter length and shorter diameter metallic CNTs 201 emerge farther out along the array and earlier in time than hybrid structures with medium length and medium diameter metallic CNTs 202 ; which, in turn, emerge farther out along the array and earlier in time than hybrid structures with longer length and longer diameter metallic CNTs 203 . [0028] Likewise, amongst semiconducting CNT based hybrid structures 204 , 205 , 206 , the hybrid structures with shorter length and shorter diameter semiconducting CNTs 204 emerge farther out along the array and earlier in time than hybrid structures with medium length and medium diameter semiconducting CNTs 205 ; which, in turn, emerge farther out along the array and earlier in time than hybrid structures with longer length and longer diameter semiconducting CNTs 206 . [0029] In this regard it is important to recognize that although only electric field array interactions have been described above (e.g., by positively or negatively charging the array's posts), other types of attractive or repulsive interactions are possible. For example, the array posts may be treated so as to promote chemical reactions with specific types of Hybrid structures. A good example is a simple addition reaction. For instance, a carbonyl group is attached to the tube and the array is coated with alcohol or amine groups that can attack the carbonyl double bond. [0030] FIGS. 3 a through 3 d show different types of array designs that may be constructed to effect separation as described above. FIG. 3 a shows an array design where, in both the x and y dimensions, all posts are aligned. FIG. 3 b shows an array design where, in both the x and y dimensions, every other post is aligned. Both the array designs of FIGS. 3 a and 3 b are apt to spread out hybrid structures in both the +x and −x directions from an entrance point at the x=0 position. [0031] The design of FIG. 3 a , however, is apt to spread the hybrid structures out farther along the x axis than the design of FIG. 3 b . Both the array designs of FIGS. 3 c and 3 d are apt to spread out hybrid structures in only the +x direction. The design of FIG. 3 c , is apt to spread the hybrid structures out farther along the x axis than the design of FIG. 3 d . Distances between array posts may vary from embodiment (e.g., ranging from microns to tens of nanometers (or less as manufacturing capabilities improve). Generally, the closer the array posts are to one another, the slower the hybrid structures will emerge from the array; and, the more spread out the hybrid structures will be when they emerge from the emitting edge of the array. [0032] FIGS. 4 a and 4 b show a system for sorting CNTs as described above. According to the design of FIGS. 4 a and 4 b , post structures for the array 403 are formed with semiconductor integrated circuit manufacturing techniques (e.g., with lithographic patterned, stacked multilayer metal and/or dielectric features) upon a silicon substrate 401 and/or using standard PDMS (polydimethyl siloxane) microfluidic structures. [0033] If electrical charge is to be applied to the posts, the posts should be made of metal that are interconnected with wiring at various locations. A lid 402 (e.g., formed with a ceramic and/or glass multi-layer structure) sits atop the array's posts. Sidewalls 404 _ 1 , 404 _ 2 (e.g., again formed with lithographic patterned, stacked multiplayer metal and/or dielectric features) have embedded channels for introducing an input fluid flow 405 containing CNTs (e.g., wrapped into hybrid structures), and providing output fluid flows 406 _ 1 , 406 _ 2 , 406 _ 3 at specific locations where specific types of CNTs are expected to emerge at specific times. [0034] The process described above can also be multiplexed to increase separation efficiency. For example, if at the output of 406 - 1 of FIG. 4B another array is added (identically designed or having more optimized geometry for finer separation) a higher percentage of targeted CNTs will emerge from the second array. [0035] In the foregoing specification, the invention has been described with reference to specific exemplary embodiments thereof. It will, however, be evident that various modifications and changes may be made thereto without departing from the broader spirit and scope of the invention as set forth in the appended claims. The specification and drawings are, accordingly, to be regarded in an illustrative rather than a restrictive sense.

A method is described that involves sorting CNT hybrid structures of differing sizes by passing the CNT hybrid structures through an arrangement of fixed structures. The sorting involves a diffusion component in which shorter CNT hybrid structures are scattered through the arrangement of fixed structures to a greater degree than at which longer CNT hybrid structures are scattered through the arrangement of fixed structures.

FIELD OF THE INVENTION [0001] The present invention relates to optical measurement devices and more particularly but not exclusively it relates to a system and apparatus adapted to measure optical properties in-situ. BACKGROUND TO THE INVENTION [0002] Optical measurement devices are used in a variety of different applications. Optical measurement devices include devices that are used to measure or otherwise determine one or more properties of light such as intensity, colour, wavelength, or other characteristics. One type of optical measurement device is an optical density sensor. One type of optical density sensor is a cell density sensor which operates by shining light through a solution to a receiver. The optical density of the solution changes the amount of absorption or scattering of the passing light. The light receiver outputs a signal dependent on the intensity of the light received, which is in turn dependent on how much scattering or absorption the solution has caused. [0003] Cell density sensors are used in biotechnology, chemical, brewing, wine making, fermentation, pharmaceutical, and other sectors of industry or research. For biotech applications, cell density sensors are ordinarily used to monitor growth of living cells in a cell culture. [0004] A disadvantage with such optical measurement systems is that they typically require an onerous process in order to be used. The process includes repeated removal of a sample of the solution at consecutive time points under sterile conditions, applying that sample to a measurement device, recording the measurement and disposing of the sample. This process increases the risk of contamination, and the loss of sample volume from the solution. [0005] It is an object of the present invention to provide an improved measurement system which overcomes or at least ameliorates some of the abovementioned disadvantage or which at least provides the public with a useful choice. Other objects of the invention may become apparent from the following description which is given by way of example only. SUMMARY OF THE INVENTION [0006] In one aspect the invention consists in a measurement device adapted for in situ light intensity sensing from within an environment comprising a housing adapted to enclose a control system and fluidly seal the control system from the environment, the housing having an outer wall and a channel fluidly connected to the environment at one or more locations, the control system comprising a controller, a light receiver component and a wireless data transmitter component, the light receiver disposed within the housing to receive light from the channel and output one or more signals indicative of light intensity, and wherein the control system is configured to receive the one or more signals indicative of light intensity from the light receiver, and output a signal indicative of light intensity to the wireless data transmitter. [0007] In one embodiment, the control system further comprises a light source disposed within the housing to define a light path that extends from the light source, through the channel to the light receiver. [0008] In one embodiment, a plurality of optical elements are disposed within the light path and are arranged to prevent light travelling substantially non parallel to the optical path. [0009] In one embodiment, a plurality of optical elements are disposed within the channel and are arranged to prevent light incident to the channel. [0010] In one embodiment, the channel extends between at least two locations on the outer wall of the housing to define a fluid flow path between the at least two locations. [0011] In one embodiment, the channel defines a substantially straight path. In another embodiment, the channel defines a curved path, including for example, a channel having an ‘S’ shape. [0012] In one embodiment, the channel meets the outer wall of the housing at an acute angle in at least one location. [0013] In one embodiment, the channel is adapted to receive a spigot containing wireless power transfer electronics. [0014] In one embodiment, the outer wall of the housing is substantially spherical or at least has a substantially circular profile. For example, the outer wall of the housing has a substantially spheroid profile, including oblate or prolate profiles. [0015] In one embodiment, the outer wall of the housing is shaped to promote mobility when immersed in an environment where optical density is to be measured. [0016] In one embodiment, the measurement device is not tethered, fixed or fastened to any one particular location within the environment. [0017] In one embodiment, the control system further comprises a temperature sensor. In various embodiments, a temperature sensor is located so as to be in contact with the environment within which the device is present, for example, a liquid suspension. In one example, a temperature sensor is located proximate the outer wall of the housing and configured to provide temperature information to the controller. In one example, a temperature sensor is located proximate the light receiver and configured to provide temperature information to the controller. In a further example, a temperature sensor is located proximate the channel. [0018] In one embodiment, the measurement device further comprises a propulsion mechanism operable to propel the device when in-situ, the controller further configured to output a signal to cause operation of the propulsion mechanism. [0019] In one embodiment, the measurement device further comprises a buoyancy mechanism operable to cause floating or sinking of the device when in-situ, the controller further configured to output a signal to cause operation of the buoyancy mechanism. [0020] In one embodiment, the controller is configured to output a signal to cause energisation of the light source. [0021] In one embodiment, the wireless data transmitter is configured to transmit data to a wireless data receiving device. [0022] In one embodiment, the control system further comprises a wireless power receiver, the receiver disposed within the housing proximate the channel so as to receive wireless power signals emitted from within the channel. [0023] In one embodiment, the channel is adapted to receive a spigot containing one or more wireless power transfer components. [0024] In one embodiment, the control system further comprises a plurality of gain setting resistors and the controller is configured to change the configuration of the resistors to affect one or more of the dynamic voltage range output from the light receiver and/or the intensity of the light source. [0025] In one embodiment, the control system further comprises a wireless power receiver, the receiver disposed within the housing proximate the exterior surface so as to receive wireless power signals emitted proximate the exterior surface. [0026] In one embodiment, the wireless power receiver is configured to provide a source of received charging power to a power source. [0027] In one embodiment, the power source is configured to provide power to the control system including one or more components of the control system. [0028] In one embodiment, the housing comprises a first and a second shell section, the first shell section having an engageable sealing surface adapted to couple with an engagable sealing surface of the second shell section, and form, when engaged, a substantially hermetic shell that encloses the control system. [0029] In one example, the housing is a substantially hermetic homogeneous shell. [0030] In one embodiment, the housing further comprises at a first aperture fluidly connected to the environment. [0031] In one embodiment, a tube is disposed with the first aperture and the inside of the tube is arranged to fluidly connect with the environment and the outside of the tube is adapted to seal to housing from the environment. [0032] In one embodiment, the housing further comprises two apertures and the tube is adapted to extend from the first aperture to the second aperture to define a fluid path through the housing. [0033] In one embodiment, the tube is disposed within the light path. [0034] In one embodiment, the tube is an optically transparent material. [0035] In one embodiment, the tube is substantially cylindrical. [0036] In one embodiment, the sealing surface of each of the first and second shell sections is threaded. [0037] In one embodiment, the sealing surface of the first and second shell sections are adapted to engage by interference fit. [0038] In one embodiment, the sealing surface of the first and second shell sections are adapted to compress about an o-ring or sealing device. [0039] In one embodiment, the first and second shell sections are adapted to be chemically or thermally bonded together. [0040] In another aspect the invention broadly consists in a measurement device adapted for in situ light intensity sensing from within an environment comprising a housing adapted to enclose a control system and fluidly seal the control system from the environment, the housing comprising: [0041] a first and a second shell section, the first shell section having an engageable sealing surface adapted to couple with an engagable sealing surface of the second shell section, and form, when engaged, a substantially hermetic shell that encloses the control system. [0042] In one embodiment, the housing further comprises at a first aperture fluidly connected to the environment. [0043] In one embodiment, a tube is disposed with the first aperture, the tube adapted to fluidly connect the environment to the inside of the housing. [0044] In one embodiment, the tube is an optically transparent material. [0045] In one embodiment, the tube is substantially cylindrical. [0046] In one embodiment, the housing has two apertures and the tube is adapted to extend from the first aperture to the second aperture to define a fluid path through the housing. [0047] In one embodiment, the sealing surface of the first and second shell sections is threaded. [0048] In one embodiment, the sealing surface of the first and second shell sections are adapted to engage by interference fit. [0049] In one embodiment, the sealing surface of the first and second shell sections are adapted to compress about an o-ring. [0050] In one embodiment, the shell is substantially spherical or at least has a substantially circular profile. [0051] In one embodiment, the shell is shaped to promote mobility when immersed in the environment. [0052] In one embodiment, the shell is not tethered, fixed or fastened to any one particular location within the environment. [0053] In one embodiment, the first and second shell sections are adapted to be chemically or thermally bonded together. [0054] In another aspect the invention broadly consists in a system comprising a sensor adapted to measure optical density from within an environment and a data processing device, the sensor comprising a housing adapted to enclose a control system and fluidly seal a control system from the environment, the housing having an outer wall and a channel fluidly connected to the environment at one or more locations, the control system comprising a controller, a light receiver component and a wireless data transmitter component, the light receiver disposed within the housing to receive light from the channel and output one or more signals indicative of light intensity, and wherein the control system is configured to receive the one or more signals indicative of light intensity from the light receiver, and output a signal indicative of light intensity to the wireless data transmitter, wherein the data processing device comprises a wireless data receiver configured to receive data transmitted by the wireless data transmitter. [0055] In another aspect the invention broadly consists in a system comprising a sensor device adapted to measure optical density from within an environment and a data processing device, the sensor device comprising a wireless data transmitter configured to wirelessly transmit a signal indicative of an optical density measurement to the data processing device, and the data processing device comprising a receiver adapted to receive a signal indicative of optical density measurements. [0056] In one embodiment, the data processing device is configured to store data received by the wireless data receiver on a storage device. [0057] In one embodiment, the sensor device has an onboard data storage unit. For example, the sensor device has an on-board data storage unit to complement or substitute for the data processing device, for example as a support system in case of network malfunction and other power outages so as to retrieve data. [0058] In one embodiment, the data processing device is configured to compute one or more statistical calculations on the stored data. [0059] In one embodiment, the data processing device is configured to determine a measure of the optical density within the channel. [0060] In one embodiment, the data processing device is configured to display a value indicative of the measure of the optical density within the channel to a display. [0061] In one embodiment, the data processing device is configured to store data indicative of the time data is received from by the wireless data receiver. [0062] In one embodiment, the data processing device is configured to interface with one or more other data processing devices. [0063] In one embodiment, the data processing device comprises a wireless data transmitter adapted to transmit a signal indicative of a measurement to be taken, and the sensor further comprises a wireless data receiver adapted to receive the signal indicative of an optical density measurement to be taken. [0064] In another aspect the invention broadly consists in a measurement system comprising an in situ light intensity sensor operable to measure optical density from within an environment and a data processing device, wherein the data processing device comprises a wireless data transmitter adapted to transmit a signal indicative of a measurement to be taken, and the sensor comprises a wireless data receiver adapted to receive the signal indicative of an optical density measurement to be taken and operate to take a measurement. [0065] In one embodiment, the device is configured to perform one or more of the following steps: store data received by wireless data receiver on a storage device, compute one or more statistical calculations on the stored data, and determine a measure of the optical intensity within the channel, output a value indicative of the measure of the optical intensity within the channel to a display, store data indicative of the time data is received from by the wireless data receiver, transmit the stored data to one or more portable computation devices, or display received data. [0066] In another aspect the invention broadly consists in a control system comprising a controller adapted for use in the measurement device and configured to output a signal to energise a light source in the measurement device, receive information indicative of an optical intensity measurement from a light receiving device, output information indicative of an optical intensity measurement to a wireless data transfer device. [0067] In one embodiment, the controller is further configured to store information indicative of an optical intensity measurement received from a light receiving device. [0068] In one embodiment, the measurement device comprises a power source and the control system is configured to measure the power from the power source. [0069] In one embodiment, the control system further comprises a plurality of gain setting resistors and the controller is configured to change the configuration of the resistors to affect one or more of the dynamic voltage range output from the light receiving device and/or the light intensity of the light source. [0070] In one embodiment, the controller is a microprocessor. [0071] In one embodiment, the control system is further configured to output a signal operable to control or at least initiate operation of one or more of a propulsion mechanism or buoyancy mechanism. [0072] In another aspect the invention broadly consists in a charging station comprising a base and a spigot and adapted to support a measurement device adapted for in-situ optical density sensing from within a fluid environment, wherein the base and/or the spigot are adapted to enclose one or more wireless power transfer components, and wherein the measurement device comprises a housing adapted to enclose one or more wireless power receiver components and the housing has a channel fluidly connected to the environment at one or more locations, and wherein the spigot is adapted to engage the channel to support the measurement device on the charging station. [0073] In another aspect the invention broadly consists in a method of measuring optical density using a device, system, or housing according to any previous statement wherein the method comprises providing the device, system, or housing, operating the controller to receiving a signal from the light receiver and outputting a signal to the wireless data transmitter. [0074] In another aspect the invention broadly consists in a method of charging the power source in the device according to any previous statement, wherein the method comprises providing the charging station, providing the measurement device, engaging the measurement device with the base and/or spigot of the charging device, and energising the one or more wireless power transfer components. [0075] In other aspects, the invention relates to a device, system, housing or station as herein described or shown in any one or more of the accompanying figures. [0076] Other aspects of the invention may become apparent from the following description which is given by way of example only and with reference to the accompanying drawings. [0077] As used herein the term “and/or” means “and” or “or”, or both. [0078] As used herein “(s)” following a noun means the plural and/or singular forms of the noun. [0079] The term “comprising” as used in this specification and claims means “consisting at least in part of”. When interpreting statements in this specification and claims which include that term, the features, prefaced by that term in each statement, all need to be present, but other features can also be present. Related terms such as “comprise” and “comprised” are to be interpreted in the same manner. [0080] As used herein, when the context allows the term “proximate” includes “at”. [0081] It is intended that reference to a range of numbers disclosed herein (for example, 1 to 10) also incorporates reference to all rational numbers within that range (for example, 1, 1.1, 2, 3, 3.9, 4, 5, 6, 6.5, 7, 8, 9 and 10) and also any range of rational numbers within that range (for example, 2 to 8, 1.5 to 5.5 and 3.1 to 4.7). [0082] The entire disclosures of all applications, patents and publications, cited above and below, if any, are hereby incorporated by reference. [0083] This invention may also be said broadly to consist in the parts, elements and features referred to or indicated in the specification of the application, individually or collectively, and any or all combinations of any two or more of said parts, elements or features, and where specific integers are mentioned herein which have known equivalents in the art to which this invention relates, such known equivalents are deemed to be incorporated herein as if individually set forth. BRIEF DESCRIPTION OF THE DRAWINGS [0084] The invention will now be described by way of example only and with reference to the drawings in which: [0085] FIG. 1 shows an embodiment of the sensor in schematic form. [0086] FIG. 2 shows a cross section of the sensor of FIG. 1 . [0087] FIG. 3 shows an overview of the sensor in situ and proximate an external control system. [0088] FIG. 4 shows a side elevation of a spherical ( FIG. 4A ) and a spheroid ( FIG. 4B ) sensor, and respective charging platforms. [0089] FIG. 5 shows a cross sectional view a spherical ( FIG. 5A ) and a spheroid ( FIG. 5B ) sensor and the charging platform of FIG. 4 . [0090] FIG. 6 shows a schematic of the components of FIG. 2 in further detail ( FIG. 6A ), and a schematic of the components of FIG. 2 with a temperature sensor ( FIG. 6B ). [0091] FIG. 7 shows a graph comparing the results obtained with a sensor and a bench top spectrophotometer as described in Example 1. [0092] FIGS. 8 to 10 show light and fluid channels arranged within the sensor. [0093] FIG. 11 shows an example of a sensor having an ‘S’ shaped fluid channel. [0094] FIG. 12 shows a spherical ( FIG. 12A-12C ) and a spheroid ( FIG. 12D-12F ) sensor and a layout of the components within the sensor. [0095] FIG. 13 shows the representative spherical ( FIG. 13A-13C ) and the spheroid ( FIG. 13D-13F ) sensors of FIG. 12 with dimensions provided. [0096] FIG. 14 shows a schematic of the components of the control system, network, and user interfaces. DETAILED DESCRIPTION OF THE INVENTION [0097] Growth in a living cell is an orderly increase in the amount of cellular components. In most living organisms, growth involves the increase in cell mass, duplication of the genetic material (DNA) followed by cell division. The division of cells increases cell number and hence the concentration of cells in a growth medium. A method of estimating cell concentration is by measuring turbidity of a suspension of cells in a liquid medium using photometry. Particle size objects, such as bacteria, suspended in a liquid scatter light that passes through the suspension. This scattering reduces the intensity of the light that is directly transmitted through the suspension. To a human eye, the suspension appears to be turbid or “cloudy”. As more light is scattered with increasing cell concentration, the reduction in light intensity can be used to measure the concentration of cells. [0098] Expressing cell growth mathematically, the intensity I of the light after it has passed through a solution or particle suspension is equal to the intensity I 0 of the incident light, multiplied by 10 −N/N10 , where N is the concentration of particles in suspension and N 10 is the concentration of particles which gives a tenfold decrease in the light intensity. [0000] I=I 0 ·10 −N/N 10 (Beer-Lambert Law) [0000] On rearranging the equation and taking the logarithm to the base 10; [0000] log I/I 0 =log 10 −N/NN 10 or −log I/I 0 =N/N 10 [0099] The term −log I/I 0 is known as absorbance or optical density (OD) of a solution or suspension. Optical density is a function of the wavelength of the light and the optical path length through the suspension. Optical density or turbidity of a suspension of cells can—after calibration—be directly converted into cell concentration. [0100] Existing probes have been designed and previously reported for large-scale fermentors or bioreactors. At the laboratory level, bacteria, yeast, fungi or mammalian cells are cultured in glass flasks and incubated at set temperatures in a shaker-incubator. Briefly, live or cryopreserved cells are inoculated into a growth medium, containing required growth supplements, inside a glass flask. This forms a broth of culture media and suspended particulates. The flasks are kept inside a temperature controlled shaking incubator to induce the cells to multiply. To monitor cell growth, aliquots of culture broth are taken manually from the flasks at regular time intervals, and measured using a spectrophotometer. The accurate monitoring of cell growth is—essential for many downstream applications and this offline measurement technique is cumbersome, time consuming, and prone to contamination and human error. [0101] Embodiments of the invention relate to a sensor that is immersible in a solution the optical density of which is desired to be measured. Embodiments of the invention also relate to a system adapted for use with the sensor. The sensor is adapted to wirelessly communicate information to an information processing system and does not require manual removal of a sample or manual use of a spectrophotometer. [0102] FIG. 1 shows a particular embodiment of the sensor 15 in schematic form. The sensor has a housing 4 which substantially encapsulates a plurality of components including a light source 2 and a light receiver 5 configured to receive light from the light source 2 . The receiver 5 can directly face the light source, or the receiver 5 can be arranged such that emitted light is adequately guided by the optical properties of nearby or incident components. In some embodiments where luminance of the solution is desired to be measured, the light source is omitted. [0103] The light source is configured to emit light when provided with an appropriate electrical stimulus. In one example, the light source is a light emitting diode (LED) or similar device which emits light when a voltage is applied. In certain embodiments the light receiver is a photodiode, phototransistor or similar device. The light receiver 5 is configured to output a signal indicative of the light intensity received. The light source 2 and receiver 5 are at least closely matched in terms of the wavelengths upon which they can efficiently transmit or receive. Further, the particular operation wavelength may be selected depending on the absorption properties of a solution desired to be tested. For example, the operation wavelength is aligned with or proximate to a peak absorption wavelength of a solution to be tested to optimise absorption efficiency and dynamic range of the measurements. [0104] The light source 2 and light receiver 5 are located a specific distance apart such that an optical path is located therebetween. A channel 3 is located within the optical path such that light emitted from the light source 2 passes through the channel and to the receiver 5 . The channel 3 has at least one opening fluidly connected with the exterior of the housing 4 such that the solution or suspension fills the channel when the sensor is immersed. The intensity of the light received by the light receiver 5 is indicative of the optical density of the substance within the optical path. [0105] In some configurations the sensor has a number of collimating devices arranged within the channel 3 and optionally also the light path between the light source 2 and light receiver 5 to lower the acceptance angle of the light reaching the receiver and to reduce the amount of light scattered, or reflected off the sides of the channels into the light receiver. Preventing or mitigating the amount of ambient or scattered light from reaching the receiver improves measurement performance. The collimating devices comprise reflective, absorptive, or dispersive optical components having a geometry that provides scattering of the light incident upon it. FIGS. 8 to 11 illustrate configurations of the sensor with ridge like collimating devices 23 arranged within the channel 3 that have the effect of enabling a lower acceptance angle of light to the light receiver, and reducing the reflectance of the channel. [0106] FIGS. 8 and 9 illustrate collimating devices 23 arranged within the light path 26 and transparent members 27 fluidly sealing the light source and light receiver portions of the light path from the channel 3 . The collimating devices 23 are optimally arranged when only light travelling substantially parallel to the channel or light path enters the light receiver. Arrangement of the light path and channel in a substantially perpendicular geometry further improves measurement performance. [0107] In some embodiments the housing has two openings such that the channel 3 extends from one side of the housing 4 to another to create a fluid flow path. This allows fluid to flow through the channel 3 as it is circulated by the natural stirring motion induced by a shaking incubator. [0108] In some embodiments the channel is a cylinder which has a curved inner surface shape that advantageously reduces the chance of bubbles forming in the cavity and affecting the light path. In other embodiments, such as those where bubbles are not a concern, the channel is circular or square or polygonal in cross section. The use of glass or other hydrophilic material to form the channel decreases the tendency of bubbles to stick to the surface. In certain embodiments the channel is sealed to the sensor shell by mechanical seals such as o-rings, or it is chemically bonded by materials such as viscous sealant. [0109] In some embodiments the channel 3 has entry ports 25 angled with respect to the sensor surface or a channel that is at least a serpentine shape. FIG. 11 shows an example of a sensor having an ‘S’ shaped channel 25 and angled entry ports. The ‘S’ shaped channel also helps to prevent ambient light from reaching the light receiver by blocking the line of sight trajectory of light entering the tunnel. [0110] In some configurations a plurality of fins 24 are arranged on the sensor surface such that when immersed within a stirred solution the fins cause the sensor to spin. The fins in combination with the ‘S’ shaped channel promote pumping of the fluid through the channel. A shorter light path through the fluid enables use in higher OD solutions and/or the use of a lower power light source. [0111] In one embodiment the sensor housing could be constructed using 3D printing techniques. A minimum wall thickness is used to ensure structural integrity is maintained during the sensor lifetime. Selective Laser Sintering 3D printing limitations for acrylic based photopolymer or nylon plastic: Objects must be manifold, minimum detail of 0.2 mm, minimum wall thickness of 0.7 mm, maximum temperature of 80° C. [0112] In various embodiments the sensor housing is formed from multiple shell components, such as two hemispheres. The shell components are, for example, chemically or mechanically fastened together to encapsulate the internal components. Mechanical fastenings to secure each shell component include threaded or interference type connections. In some embodiments the shell components comprise two substantially hemispherical shell sections. [0113] Each shell section has a mating surface where the shell sections are to oppose and engage. For example, the mating surface has a threaded connection complimentary to the opposing shell sections such that the shell sections can be screwed together. Alternatively, the mating surface of each shell section is sized to engage with an interference fit and therefore allow the use of non circular engaging surfaces. The housing is then completed by applying pressure to join the two shell sections and force the engagement of the mating surfaces. One or more sealing devices such as o-rings or semi viscous sealant may be employed to ensure leak proof engagement of the shell sections. In another embodiment, sealing is by friction, welding, or other engagement to ensure leak proof engagement, for example so as to form a homogeneous surface. [0114] An optimum form of the sensor 15 for facilitating movement when immersed in a solution is that of a spherical form such as shown in the figures. However, in certain embodiments other forms that facilitate movement of the sensor, or at least do not substantially prevent movement, within a moving solution are utilised. For example, a substantially cylindrical or elliptically shaped housing form, a spheroidal form, or any other curved surface shaped housing form is utilised. [0115] The outer surface of the sensor can include dimples or spikes. Impressions can help reduce fluid drag or provide traction to resist or reduce movement in fast moving fluids. Spikes can be provided where it is desired that the sensor embed itself in material on the bottom of the solution. For example, the sensor can be used in a stream of fluid where measurements are desired to be taken, such as a riverbed. The spikes help to fix the sensor in one location with respect to the stream. [0116] To further facilitate movement of the sensor within an agitated solution, the outer surface of the sensor may include one or more fluid dynamic surfaces operable to impart kinetic energy to the sensor from the solution, or by an induced rotation of the shell with respect to the solution. The surfaces may comprise fins, contours, impressions or depressions joined to or formed in the outer sensor surface. [0117] In some embodiments the sensor is constructed to have an eccentric weight distribution. This can assist in movement of the circulating fluid through the channel, or to preferentially align the channel with a particular direction, such as that of a flow path of a solution. [0118] In some embodiments the light source 2 and receiver 5 are mounted to an electronics substrate 1 that is, in turn, releasably mounted to or within the housing 4 . When mounted to the substrate 1 , the light source and receiver may be readily removed from the sensor 15 and replaced with other combinations. This may allow the selection of particular operation wavelengths and light intensities depending on the solution desired to be tested. In some embodiments the substrate 1 is a circuit board able to flex such that it may easily fit within packaging constraints that sensor housing may impose. However, it is ideal that the light source 2 and receiver 5 are rigidly spaced apart such that alignment is maintained and movement or vibration has no substantial effect on the accuracy of the optical components. One or more guide lugs can be provided within the housing to facilitate repeatable and stable mounting of the substrate 1 . [0119] FIG. 2 shows a cross section of the sensor of FIG. 1 showing the channel 3 disposed in the light path 9 between the light source 2 and the light receiver 5 . The channel 3 is open to receiving solution 8 which causes absorption or scattering of light within the channel and light path. [0120] FIG. 2 also shows a control unit or controller 10 that is configured to connect to the light source 2 , the light receiver 5 , a power source 6 and a wireless communications interface 11 . The control unit is a microprocessor having at least one or more analogue-to-digital converter (ADC) inputs, one or more digital outputs and/or serial data transmission and receiver pins for communication with external protocol capable devices. [0121] FIG. 6 shows a schematic of the components of FIG. 2 in further detail ( FIG. 6A ), and a schematic of the components of FIG. 2 with a temperature sensor 60 ( FIG. 6B ), and in particular shows the controller 10 and connected components. The controller is configured to connect with the wireless communications interface 11 to at least transmit and also receive data from an external system. [0122] In some embodiments the wireless communications interface 11 is a radio transceiver. However, in other embodiments the interface 11 may be a transmitter only. The interface can be, for example, 2.4 GHz transceiver having a common communications protocol such as a Bluetooth transceiver. Other transmission frequencies and protocols may be used. In circumstances where long range communication is required, or communication through matter having substantial radio frequency attenuation, it can be beneficial to use lower frequencies such that wirelessly transmitted power can be kept relatively low to conserve battery power. However, higher transmission frequencies offer benefits such as smaller antennas and may therefore be most appropriately applied in circumstances where limited packaging space is available within the sensor housing 4 . Further, in some embodiments the PCB track or flexible antennas are incorporated. Alternatively, the antenna extends externally to the sensor housing as long as appropriate shield materials suitable for sterilisation are used. [0123] In various embodiments the controller 10 is configured to control energisation of the light source 2 and receive a signal from the receiver 5 indicative of the amount of light received. The controller can facilitate an automated process where, for example, the light source 2 is periodically energised and the receiver 5 output received and stored. Alternatively, the controller can respond to an instruction received via the communications interface 11 to make a measurement. [0124] The controller 10 samples a voltage received from the light receiver 5 via an ADC input pin at periodic intervals. The signal the controller receives from the light receiver 5 is indicative of the light scattered or absorbed by solution within the optical path 3 . The optical density of the solution can be determined from the received signal. As the solution becomes more optically dense, the intensity of the light received by the light receiver is reduced. [0125] In some embodiments the controller has an ADC configured to sample the signal received from the light receiver. The sampled value can optionally be converted a measurement via Beer's Law or stored for later use and/or transmission to an external system. The controller can process the measurement internally, for example using preconfigured software, or the controller can output the raw value to the communications interface and an external processing system may then calculate the measurement value. The controller has a digital output configured to control energisation of output of the light source 2 . The digital output may be configured to provide, for example, a PWM signal representing a desired intensity output. Such a PWM signal may be amplified by appropriate electronics should the controller output not be able to supply enough current on its own. [0126] FIG. 14 presents a schematic depicting one embodiment of the invention, in which one or more sensor devices 15 are controlled by a control system 141 (typically a PC or laptop) in the proximity of the one or more sensor devices. The control system is connected 143 to the one or more sensor devices, for example via Bluetooth 142 , and to a computer network 144 for user interaction. Users 145 can remotely control or monitor the sensor devices by connecting to the control system. In one example, one or more of the sensor devices have an identifier 146 , such as an LED, to differentiate the devices present. The identification process and other sensing aspects of the device can be triggered via software 147 running on the control system. In various embodiments the software is configured to identify one or more of the sensor devices, for example internally by a unique ID, to make it identify itself to the user (for example by LED), to control the measuring functions of the sensor device, to fetch recorded data from the sensor device, and/or to display and/or analyse the data. [0127] To minimise power consumption, in some embodiments the controller is configured to pulse the light source or energise it only for short periods such as when the receiver output is being monitored. Energising the light source for at least 10 ms ensures that light source avoids detecting any transients and that the receiver output is likely to be stable. The light receiver receives enough light in order for it to make correct measurements and that intensity of the light source is constant. [0128] The light receiver output is tuned such that it remains within a linear range and does not saturate. The linearity of response is ensured either by selection of the components during construction or dynamically by a configuration of programmable resistors that are connected to the controller to form voltage divider circuits and/or control the gain of an active signal amplifier. The programmable resistors can be set based on knowledge of the dynamic output range of the light receiver and/or light source to tune either the output sensitivity or intensity respectively. This allows the sensor to be configured to measure a wide range of optical densities and that configuration is changeable using the controller to implement changes in sensitivity and intensity. [0129] In some embodiments, the sensor may incorporate a propulsion mechanism operable to provide motility of the sensor within a vessel. For example, in environments with large fluid volumes, the propulsion mechanism advantageously enables the sensor to operate to sample from several locations within the vessel, and wirelessly transmit the sample to a distal location. Propulsion can be achieved, for example, by having a rotatable fin mounted external to the sensor housing. Rotating the fin by way of a motor propels the sensor within the solution. The propulsion mechanism can also be used to replace a laboratory shaking or stirring platform by actively agitating the solution by the sensor moving in the solution and/or the sensor moving the solution relative to its position. This may be advantageous in circumstances where the optical density of a solution is desired to be known in a non-laboratory environment. [0130] In some embodiments, the sensor incorporates a buoyancy control device. For example, buoyancy control may be desired in environments with vertically large fluid volumes such as beverage fermentation vats. The buoyancy control device advantageously enables the sensor to take measurements from many vertical locations as it rises and sinks. Alternatively, the buoyancy of the sensor could be selected to float or sink in a particular solution to be tested. Buoyancy control can be achieved, for example, by compressing a compressible fluid with a piston to change the internal density of the sensor. Alternatively, a fluid bladder can be used to draw solution into the bladder to change the buoyancy. [0131] In certain embodiments, particularly where the sensor includes a buoyancy and/or a propulsion mechanism, the controller can be configured to actively control propulsion and/or buoyancy of the sensor in-situ. For example, the controller is configured to have the sensor move in the solution while recording measurements. For example, the sensor may be located in a vessel having a large vertical distance such as a fermenter. The sensor can travel the vertical distance by control of buoyancy and/or propulsion while also recording measurements to attain a continuous profile of the vessel. In other embodiments the buoyancy and/or a propulsion mechanism is operated in free form or a predetermined activation pattern. For example, when the buoyancy or propulsion mechanism is configured to cause the sensor to rise and/or sink one or more times. [0132] In some embodiments the sensor includes a second light receiver and optical components configured to reflect a portion of the light transmitted by the light source to that second light receiver. The light received by the second light receiver is indicative of the output power from the light source and can be used as a calibration measure. [0133] In some embodiments the sensor includes one or more temperature sensors arranged within the housing and configured to provide temperature information to the controller. For example, temperature sensors located proximate to the light receiver can be used to compensate for temperature related drift of the light receiver. A temperature sensor located near the power source can be used to indicate excessive temperature generation. A temperature sensor located in contact with or proximate a surface in contact with the environment, for example near the housing surface, or proximate the channel, can be used to indicate environment temperature information, for example solution temperature information. [0134] The power source 6 is configured to provide power to the controller 10 and other components located within the housing 4 . The most useful power source is a rechargeable battery. In this configuration, a charging system 7 is connected to the battery to provide a source of power from which the battery can be recharged. In this configuration the charging system has a wireless power transfer receiver. In other configurations, the charging system and battery may be replaced with a wireless power receiving device that continuously receives power to operate the sensor, or receives power at least when measurements are desired to be taken. Inductive power transfer technology may be used to apply power to the sensor or battery charging system from a remote location without the use of a wired connection. [0135] The power system 7 incorporates appropriate electronics configured to convert a received wireless power signal into a voltage useful for charging the battery or powering electronics within the housing. The particular electronics and configuration required are dependent upon incoming wireless power transfer signals and the particular devices to be powered. Those skilled in the art will recognise the need for the electronic circuits to be tailored to the requirements. However, it is noted that rectification and/or DC to DC conversion circuits are most applicable. The controller 10 may further facilitate power management by, for example, monitoring the voltage of the battery 6 and communicating readings either facilitating transmission of a signal indicative of the need for recharging or automatically activate a recharging process. [0136] In some embodiments the wireless charging system has a coil operating with a voltage of around 5V and frequency between 112-205 kHz on a 100 kHz tuned coil circuit with 5 W max power output. For optimum charging operation, the charging components can include Qi compliant inductive charging with device detection, power transmission management and foreign object detection. [0137] To optimise sensor operation, the power management electronics and the communication electronics are physically separated by a practical distance to mitigate or eliminate electromagnetic or radio frequency interference creating undue noise. Further, separate ground planes between power management electronics and the communication electronics is beneficial to further isolate noise. [0138] To provide power to the sensor a battery can be a single cell 3.7 V lithium ion polymer 110 mA/hr battery having 200 mA discharge and 100 mA charge rates. However any high energy density rechargeable battery could be employed. Alternatively, where the size of the sensor housing is not limited, lower density energy sources could also be used. [0139] In some embodiments, one or more electromechanical kinetic energy harvesting mechanisms may be used in place of, or alongside, a wireless power transfer device to facilitate a source of power to recharge the battery 6 . In this way, movement of the sensor while in use can generate electrical energy used to power onboard electronics or charge a battery. The time required to charge the battery in the sensor may be therefore reduced or not required. [0140] The controller 10 may further be configured to use the light source 2 as a status indicator in configurations where the light source can be seen from outside the vessel containing the sensor. For example, the light source can be flashed to show a code indicative of parameters to a user in visual range. The code can be indicative of information such as, full memory or low battery, or for identification of a particular sensor in an environment where many sensing devices are simultaneously deployed. [0141] FIG. 4 shows a side elevation of a side elevation of a spherical ( FIG. 4A ) and a spheroid ( FIG. 4B ) sensor 15 residing on a spigot 21 that forms part of a charging platform 22 . FIG. 5 shows a cross sectional view the sensor 15 and charging platform 22 of FIG. 4 , and a perspective view of a charging platform 22 without a sensor attached. While the channel 3 of the sensor 15 performs the task of providing an opening that allows the solution to flow into the optical path between the light transmitter 5 and light receiver 4 , it also provides a mounting receptacle that allows the sensor to be mounted securely to the platform. The spigot 21 of the platform 22 incorporates a wireless power transmission, or inductive charging device. The sensor 15 has the wireless power receiving device positioned proximate to the spigot 21 when mounted on the spigot to optimise wireless power transfer efficiency. [0142] When the sensor 15 is not in use, it can be placed upon the platform 22 which then provides wireless power to charge the built in battery 6 . The sensor 15 does not therefore require a wired interface for power transmission or recharging and the sensor has a charged battery when it is required to be used. [0143] FIG. 3 shows an overview of the sensor 15 in-situ and proximate an external control system. The external control system is configured to work harmoniously with one or more sensors by being configured to respond to communication signals transmitted from one or more sensors, store data received from the one or more sensors and optionally display data. [0144] The external control system comprises one or more computational devices and may comprise one or more of a stand-alone computer, laptop 20 , smart phone 19 or tablet type device. A base station 18 may optionally be provided to interface one or more computational devices to one or more sensor devices 15 . The base station 18 may comprise, for example, a wireless communications interface complementary to the wireless communications interface incorporated in the sensor. [0145] The base station 18 may also comprise a wireless power transfer device adapted to provide a wireless power transfer signal to the sensor. The base station 18 may further comprise computational ability and provide a replacement for other computational devices. The base station 18 may further comprise one or more display devices adapted to display data such as a real time optical density measurement or sensor battery capacity status. [0146] The sensor 15 is shown in use whereby it is immersed within a vessel 13 also containing a solution from which an optical parameter, such as the optical density, is desired to be measured. The vessel could be a beaker, flask or similar container. The vessel 13 is optionally located within an incubation cubicle 12 for control of the environmental temperature. The vessel optionally resides atop a plate that provides mechanical movement to the vessel to simulate stirring or agitation of the solution within. [0147] Cell growth rate and/or density in a solution containing live cells can be determined by periodic measurement of optical density of the solution within the sensor channel 3 . The sensor 15 determines a measure indicative of the optical density of the solution 8 and wirelessly transmits a signal indicative of the optical density to the base station 17 for further interpretation. The wirelessly transmitted data may optionally include identification information in the event several sensors are used in close proximity. In this way, the base station may determine the particular sensor from which a signal was received. [0148] The base station may also be configured to transmit a signal to the sensor and the sensor configured to receive that signal and respond appropriately. For example, the base station may be configured to transmit a signal indicative of a sample value to be taken by the sensor. The sensor is configured to receive that signal, determine a measurement from the solution in the channel 3 and transmit data back to the base station 17 . Alternatively, the sensor may transmit blocks of data at periodic intervals which allows the sensor to sample over a longer time period than continuously transmitting data. This minimises the energy consumption associated with data transmission, or for extending operation time should the battery energy be depleted. The base station may further be configured to control activation of any propulsion or buoyancy control within the sensor. [0149] FIG. 12 shows a spherical embodiment ( FIG. 12A-12C ) and a spheroid embodiment ( FIG. 12D-12F ) of the sensor 15 optimised for use in a laboratory vessel and with an optimised layout of the internal components. In particular, FIG. 12A shows a front view and cross sectional view AA, FIG. 12B shows a side view and cross sectional view BB, and FIG. 12C shows a top view and cross sectional view CC. [0150] The sensor 15 has two joinable hemispherical shell sections, an upper shell section 31 and a lower shell section 32 . The shell sections can be joined by a releasable mechanism such as an interference fit or threaded connection. [0151] The channel 3 is formed from a borosilicate (for example, Pyrex™) tube that extends from one extent of the shell to the other. The channel 3 is sealed to the shell sections by o-rings 30 . A light source 2 and light receiver 5 are disposed about the channel. [0152] A collimator 23 is provided proximate the light source and receiver to guide light through the channel, minimise light scattering and minimise ambient light from entering the light receiver. [0153] A battery 6 is located in one portion of the housing and one that is distant from the location of the wireless communication device 11 to minimise shadowing of radio signals. The wireless communication device 11 is a Bluetooth transceiver module. A controller 10 is connected to the wireless communication device 11 , the light source 2 and the light receiver 5 . [0154] FIG. 12D shows a front view and cross sectional view AA, FIG. 12B shows a side view and cross sectional view BB, and FIG. 12C shows a top view and cross sectional view CC, of the spheroid sensor 15 . [0155] The sensor 15 has two joinable hemispherical shell sections, an upper shell section 31 and a lower shell section 32 . The shell sections can be joined by a releasable mechanism such as an interference fit or threaded connection. [0156] The channel 3 is formed from a borosilicate (for example, Pyrex™) tube that extends from one extent of the shell to the other. The channel 3 is sealed to the shell sections by o-rings 30 . A light source 2 and light receiver 5 are disposed about the channel. [0157] A collimator 23 is provided proximate the light source and receiver to guide light through the channel, minimise light scattering and minimise ambient light from entering the light receiver. [0158] A battery 6 is located in one portion of the housing and one that is distant from the location of the wireless communication device 11 to minimise shadowing of radio signals. The battery is connected to a charging coil or coils 121 in the sensor 15 , in which current is induced when the sensor 15 is placed on the charging platform 22 by a charging coil or coils 121 , optionally forming part of a charging PCB 122 , present in the charging platform 22 and the charging platform is powered, for example via a DC jack 123 . The wireless communication device 11 is a Bluetooth transceiver module. A controller 10 is connected to the wireless communication device 11 , the light source 2 and the light receiver 5 . [0159] FIG. 13 shows the representative spherical ( FIG. 13A-13C ) and the spheroid ( FIG. 13D-13F ) sensors of FIG. 12 with dimensions provided. In particular, FIG. 13A shows a cross sectional view AA as depicted in FIG. 13B , FIG. 13C shows a cross sectional view BB as depicted in FIG. 13B , FIG. 13D shows a cross sectional view GG as depicted in FIG. 13E , and FIG. 13F shows a cross sectional view EE as depicted in FIG. 13E . [0160] Key dimensions in the representative spherical sensor of FIG. 13A-13C include the outer sensor dimension of 40 mm and the channel 3 inner dimension of 7 mm. The width of the beam of light is 1 mm and the distance between the light source and detector is 22 mm. Key dimensions in the representative spheroid sensor of FIG. 13D-13CF include the outer sensor diameter of 33 mm and length of 39 mm, and the channel 3 inner dimension of 4 mm. The width of the beam of light is 1 mm and the distance between the light source and detector is 13 mm. [0161] Use of the sensor 15 may further include one or more of the following steps, in any order: The sensor is placed on the charging base such that the spigot 21 is inserted within the channel 3 that extends through or partly through the sensor housing. The sensor is charged using the wirelessly coupled base station 22 or similarly capable device. The transmitter in the sensor is connected with a receiver to facilitate wireless communication therebetween. The sensor is placed inside a vessel containing a solution to be measured. The vessel containing the sensor is placed upon a mechanical stirring or shaking device such as a shaker incubator to agitate the solution in the vessel. The base station transmits a signal to the sensor when a measurement is desired. The controller receives a signal that a measurement is desired, causes energisation of the light source, measures the signal from the light receiver and outputs information to be transmitted to the base station. The controller is configured to cause energisation of the light source, measure the signal from the light receiver and output information to be transmitted to the base station. The controller is configured to cause energisation of the light source, measure the signal from the light receiver and store measured signals to later be transmitted to the base station. The controller measures the remaining battery capacity. If the capacity is determined to be low, an additional energy harvesting mechanism such as an electro-mechanical generator may operate to convert kinetic energy of the sensor into electricity to charge or supplement the battery. Cleaning the sensor using, for example, a 70% ethanol solution, for re-use. [0173] The sensor may operate to achieve a measurement by using one or more of the following steps, in any order: Energise the light source 2 for a predetermined period of time. For example, the light source is energised for approximately ten milliseconds. Measure one or more samples indicative of the light received by the light receiver 5 . For example, approximately five samples are recorded by the controller. Perform a statistical calculation of the measured samples. The most useful statistical calculation is where the average of several samples is calculated. Those skilled in the art will appreciate that other statistical or filtering calculations could be performed, or several performed and combined where circumstances dictate that this would provide a more meaningful measure. [0177] The sensor has numerous advantages, including: [0178] A housing that entirely seals internal components from the environment whilst being able to be submerged within a solution to be tested. The sensor can be constructed from a material that is easily chemically sterilised for repeated use in a variety of applications. The sensor is sealed and self contained which helps to prevent contamination. The sensor is stored and used remotely from the base station without requiring a wired interface. The sensor can be transported with a vessel containing a solution, for example, between a storage area and a shaking table where a solution is to develop. The sensor is able to be mobile while immersed in a solution thereby improving the ability of the sensor to provide measurements from a variety of locations within the vessel. EXAMPLE 1 Assaying Optical Density of Microorganism Culture [0184] Preliminary testing of a sensor 15 was conducted alongside measurements gathered using an Eppendorf BioPhotometer Plus bench top spectrophotometer simulating a typical yeast growth assay. [0185] A widely used laboratory strain Saccharomyces cerevisiae: W303 was cultured overnight to exponential phase in YPD complete media and the cells were collected by centrifuge. [0186] The sensor 15 was sterilized in 70% ethanol and added to fresh media in a sterile culture flask before being placed inside an incubator containing an orbital shaker. [0187] A baseline measurement was taken following a period of approximately ten minutes to allow for temperature equilibration. The base line data output was stable regardless of shaker motion, internal lighting, or shielding of external ambient light. Cells were added to the media at an initial inoculation to OD 600nm 0.1 as measured by the Eppendorf spectrophotometer. [0188] Cells were then added to the culture media at a volume equivalent to 0.1 OD 600nm units and at each point measurements were made in parallel using both the commercial spectrophotometer and sensor 15 up to a final density of 0.8 OD 600nm . [0189] As measurements using the commercial spectrophotometer involved removal of a small sample from the culture, dilution and pipetting into a cuvette, there is some error resulting from the numerous small volume liquid handling steps which is apparent in the data supplied and typical for this type of measurement. [0190] FIG. 7 shows a comparison of the sensor 15 and the bench top spectrophotometer. In-situ photometric device measurements were obtained by averaging 14 samples each consisting of an averaged value of 5 intensity measurements. Manual triggering (in rapid succession) of the measurements were used to obtain the intensity data after addition of each aliquot of cells. [0191] As can be seen, the data provided by the in-situ device of the present invention provides a response to the change in optical density during the course of the experiment. In the case of non-linear response due to saturation, the non-linearity can be corrected by precalibration of the device over the expected optical density range, before deployment. In the case of remaining in the linear regime, no such calibration is required and the results may be directly interpreted and in this case accuracy at least comparable to standard measurement equipment is obtained. [0192] In this specification, where reference has been made to external sources of information, including patent specifications and other documents, this is generally for the purpose of providing a context for discussing the features of the present invention. Unless stated otherwise, reference to such sources of information is not to be construed, in any jurisdiction, as an admission that such sources of information are prior art or form part of the common general knowledge in the art. [0193] Where in the foregoing description reference has been made to elements or integers having known equivalents, then such equivalents are included as if they were individually set forth. [0194] Although the invention has been described by way of example and with reference to particular embodiments, it is to be understood that modifications and/or improvements may be made without departing from the scope or spirit of the invention. INDUSTRIAL APPLICABILITY [0195] The devices and systems of the invention, and the methods of using them have application in a wide range of industries and environments, including medical, biotechnological and pharmaceutical research and production, food and beverage technologies, industrial processing, the horticultural and agricultural sectors, and others.

The present invention relates to optical measurement devices and systems, and methods of using these systems and devices, and more particularly but not exclusively it relates to a system and apparatus adapted to measure optical properties in-situ.

BACKGROUND OF THE INVENTION The present invention relates generally to a striker for a vehicle closure, and in particular to an adjustable striker assembly for engaging a latch on a compartment door. Closures in vehicles commonly have a latch and striker type of arrangement for holding the closure in its closed position. Often, the latch will be mounted to the closure with a button or other type of release mechanism controlling the latch. A striker is then mounted to a vehicle component or structure and located so the latch will engage the striker when the closure is moved to its fully closed position. The latch engagement with the striker then holds the closure in its closed position until the button is actuated to release the latch from the striker. A common use for a latch and striker assembly is a vehicle glove box, with the latch mounted to the glove box door and the striker mounted to the glove box portion of an instrument panel. The glove box door is fitted to provide a flush fitting condition. In modern vehicles, the styled surfaces of instrument panels typically do not provide natural overhangs or features that will hide improper door fit. Consequently, for proper aesthetics, automotive glove box doors are required to fit nearly perfectly on every vehicle. When an automotive instrument panel is assembled and the various components are fitted, the flushness can be accomplished relatively easily, resulting in good appearance and satisfying the desired styling look for the instrument panel assembly. After this instrument panel assembly is shipped from the supplier to the automotive assembly plant and attached to the vehicle body, however, these fits invariably change. This may be due to changes in the instrument panel assembly during shipping and handling, as well as variances in each automotive body that may distort the instrument panel assembly when attached to the automotive body. The change in the fit of the glove box door, then, may necessitate an adjustment to assure the latch and striker assembly engage properly. Typical strikers used with glove box closures are made of bent steel wire, which is welded to a steel plate that is then riveted or screwed to the instrument panel. Many times the strikers are not precisely located when installed on the instrument panel—due to tolerances in locating the plate or twisting that may occur when mounting screws are tightened. These tolerances may add to the concern with the glove box door fit. Because the flush fit of the door is important to the appearance and quality of the vehicle interior—despite the fact that every instrument panel is installed into a vehicle body that has some variation due to manufacturing tolerances—the strikers must be readjusted by assembly personnel on most vehicles. In addition, the need arises—whether due to instrument panel installation variations, striker installation variations, or both—to adjust the striker so the latch will engage properly with it when the glove box door is closed. Conventionally, this adjustment is accomplished by guessing what adjustment is needed and manually bending the wire of the striker. This is a very crude and inexact process that often results in the striker being bent and moved in unintended directions. The unwanted distortion from this crude adjustment process may increase the friction between the latch and striker and so may raise operating efforts—even possibly cause some binding in the latching assembly. Moreover, this crude alignment process may add to labor costs, and also may risk damaging the striker. Thus, the adjustment process, while assuring that the latch will engage the striker, may prevent the smooth operation of the latch and striker assembly and add to the vehicle assembly costs. It is desirable, therefore, to provide a striker that is used with a latch on a vehicle closure that allows for easy and accurate adjustment of the striker to assure that the latch and striker assembly works smoothly and properly when the closure is properly aligned with its compartment. SUMMARY OF THE INVENTION An embodiment contemplates a striker assembly for use with a closure for a vehicle compartment. The striker assembly may comprise a striker and a bracket. The striker may have a main shaft with a first end and an opposed second end, a hoop extending from the first end, threads extending around the main shaft, and a rotational retention feature adjacent to the threads on the main shaft. The bracket may have a pair of mounting flanges adapted to be mounted to one of the closure and the vehicle compartment, a first wall including a threaded hole therethrough that operatively engages the threads on the main shaft, a second wall, spaced from the first wall, having a hole therethrough that engages the main shaft to allow for axial and rotational movement of the main shaft relative to the second wall, and a flexible arm having a catch operatively engageable with the rotational retention feature such that the flexible arm is in a substantially relaxed position when the catch is aligned with the rotational retention feature and is in a flexed position when the catch is not aligned with the rotational retention feature. An embodiment contemplates a method of forming an adjustable striker assembly that can be used with a closure for a vehicle compartment, the method comprising the steps of: forming a striker having a main shaft with a first end and an opposed second end, a hoop extending from the first end, threads extending around the main shaft, and a rotational retention feature adjacent to the threads on the main shaft; and insert molding a single piece bracket around the striker, including forming a first wall including a threaded hole therethrough that engages the threads on the main shaft, a second wall, spaced from the first wall, having a hole therethrough that engages the main shaft, and a flexible arm having a catch engaging the rotational retention feature such that the flexible arm is in an as molded position when the catch is aligned with the rotational retention feature and is in a flexed position when the catch is not aligned with the rotational retention feature. An embodiment contemplates a striker assembly mounted to a vehicle glove box and adapted to engage a latch of a glove box door. The striker assembly may comprise a striker having a main shaft with a first end and an opposed second end, a hoop extending from the first end and adapted to engage the latch, threads extending around the main shaft adjacent to the second end, and a rotational retention feature on the main shaft between the threads and the hoop; and a bracket having a pair of mounting flanges mounted to the vehicle glove box, a first wall including a threaded hole therethrough that operatively engages the threads on the main shaft, a second wall, spaced from the first wall, having a hole therethrough that engages the main shaft adjacent to the first end, and a flexible arm having a catch operatively engageable with the rotational retention feature, the bracket being a single, monolithic piece. An advantage of an embodiment is that the adjustable striker assembly is relatively quick and easy to adjust to assure proper alignment with a corresponding latch. In addition, the initial torque required to initiate the adjustment is relatively low because there need only be one V-groove engaged with one flexible arm to hold the striker in the correct position after adjustment. And the adjustment friction (torque required to rotate the striker and effect the adjustment) is relatively low due to the fact that there only needs to be one engaged adjustment thread between the bracket and striker. An advantage of an embodiment is that the adjustable striker assembly can be adjusted while minimizing the risk of damaging the striker, thus assuring smooth operation of the latch and striker assembly without increased friction or binding concerns. An advantage of an embodiment is that the adjustable striker assembly may not require assembly of separate components since the bracket body can be formed on the striker. With only two parts, the striker assembly is relatively simple and durable. An advantage of an embodiment is that the adjustable striker assembly can be adjusted easily without tools, and yet maintain the proper position after adjustment. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is a perspective view of a striker assembly according to a first embodiment. FIG. 2 is a side view of the striker assembly of FIG. 1 . FIG. 3 is an end view of the striker assembly of FIG. 1 . FIG. 4 is a plan view of the striker assembly of FIG. 1 . FIG. 5 is a perspective view of a striker assembly and mounting location according to a second embodiment. FIG. 6 is a plan view of a striker assembly according to a third embodiment. FIG. 7 is a section view taken along line 7 - 7 in FIG. 6 . FIG. 8 is a section view taken along line 8 - 8 in FIG. 6 . FIG. 9 is a partial section view taken along line 9 - 9 in FIG. 6 . DETAILED DESCRIPTION FIGS. 1-4 illustrate a striker assembly, indicated generally at 20 , which engages with a latch (not shown) of a latch and striker assembly, indicated generally at 22 . The striker assembly 20 is mounted to one of a door or compartment (not shown in the first embodiment), with the latch being mounted to the other of the door or compartment. The striker assembly 20 may be secured to the door or compartment with rivets (not shown in the first embodiment) or other suitable fastening mechanisms. The compartment may be, for example, a glove box opening in an instrument panel. The striker assembly 20 includes a striker 30 . The striker 30 may be made of, for example, a metal such as steel—although, other suitable materials may be employed instead. The striker 30 includes a main shaft 32 , which is generally cylindrical, having a first end 34 from which a D-shaped hoop 36 extends, and a second end 38 from which a retention flange 40 extends. Both the retention flange 40 and the hoop 36 will resist impact loading, preventing the striker 30 from being either pulled or pushed out of a bracket (discussed below) of the striker assembly 20 . The hoop 36 may be formed into shapes other than the D-shape, if so desired. The striker 30 also has threads 42 formed onto its main shaft 32 adjacent to the retention flange 40 . Recessed within the main shaft 32 , between the threads 42 and D-shaped hoop 36 , are two V-grooves 48 . These V-grooves 48 are oriented to extend longitudinally on the main shaft 32 and be located one hundred eighty degrees apart from each other. While having the V-grooves 48 recessed into the main shaft 32 creates one type of catch, other types of rotational retention features that will selectively catch and release the striker 30 from the bracket (discussed below) for selective rotation may be employed instead, if so desired. Moreover, while it is preferred to have a pair of V-grooves 48 , the striker 30 may include only one if so desired. The striker assembly 20 also includes a bracket 60 . The bracket 60 may be formed from, for example, molded plastic, or other suitable materials. The bracket 60 has mounting flanges 62 with mounting holes 64 for receiving the fasteners when securing the striker assembly 20 to the compartment or door. Side support walls 65 —extending generally parallel to the main shaft 32 and normal to the plane of the mounting flanges 62 —define a central opening 66 through which the striker 30 extends. Rear support walls 74 extend from the side support walls 65 , adjacent to and on both sides of the retention flange. The bracket 60 includes a trunnion plate 68 that surrounds and supports the main shaft 32 of the striker 30 at an end of the central opening 66 adjacent to the D-shaped hoop 36 . The trunnion plate 68 defines an opening 69 that is generally smooth and circular where it contacts the main shaft 32 of the striker 30 . Thus, when the striker 30 is rotated relative to the bracket 60 , the trunnion plate 68 allows rotation without significant resistance, while providing support for the striker 30 (i.e., providing a bearing function). The bracket 60 also includes an end wall 70 that surrounds and supports the main shaft 32 at the location of the threads 42 , adjacent to the retention flange 40 . The end wall 70 extends parallel to the trunnion plate 68 . The end wall 70 defines a threaded hole 72 that engages the threads 42 on the main shaft 32 . Accordingly, when the striker 30 is rotated relative to the bracket 60 , the engagement of the threads 42 with the threaded hole 72 will cause the striker 30 to also move axially relative to bracket 60 by an amount determined by the pitch of the threads 42 . The threaded hole 72 also provides a bearing surface for the striker 30 . A flexible arm 76 extends from one of the side support walls 65 into the central opening 66 , and has a free end 80 adjacent to the main shaft 32 of the striker 30 . Near the free end 80 , a catch, such as a barb 78 , extends from the arm 76 toward the main shaft 32 and is shaped to engage the V-groove 48 . The axial location of the barb 78 relative to the V-groove is such that, as one rotates the striker 30 relative to the bracket 60 , the barb 78 will engage each V-groove (when rotationally aligned) for the full distance of axial travel (in both directions) of the striker 30 relative to the bracket 60 . The barb 78 and V-grooves 48 are located so that the D-shaped hoop 36 is held in the proper rotational orientation each time the barb 78 engages one of the V-grooves 48 . In this first embodiment, the number of threads 42 , thread pitch, and length of V-grooves 48 are configured to allow for two half-turns in either direction from a nominal central position. Of course, the number of threads 42 , thread pitch, length of V-grooves 48 , and length of the main shaft 32 can be modified, if so desired, to provide for a greater or lesser amount of adjustment capability. Moreover, if so desired, the bracket may have a second flexible arm and barb (not shown), located and oriented one hundred eighty degrees from the first so that both V-grooves 48 are engaged at the same time. This may, however, add more cost and complexity and increase the adjustment torque more than is desirable. The manufacturing of the striker assembly 20 may include the bracket 60 being formed around the striker 30 , if so desired. For example, if the striker 30 is metal and the bracket 60 is made of plastic or some other similar, suitable material, then the striker 30 may be insert-molded into the bracket 60 . The trunnion plate 68 is molded around the main shaft 32 to form the opening 69 , and the end wall 70 is molded around the threads 42 to form the threaded hole 72 . In addition, the flexible arm 76 and barb 78 are molded so that the barb 78 is molded into one of the V-grooves 48 . Thus, when the striker 30 is in the as-molded position, the flexible arm 76 is also in its as-molded position, and so is not flexed. The flexible arm 76 is flexed when one rotates the striker 30 sufficiently to cause the barb 78 to be pushed out of the V-groove 48 . This particular manufacturing process is advantageous in that it creates a two piece adjustable striker assembly 20 that does not require assembly after forming, while still allowing for adjustment between the D-shaped hoop 26 that engages the latch (not shown) and the mounting holes 64 where the striker assembly 20 is mounted to the door or compartment (not shown in this embodiment). Also, by forming the bracket 60 this way, the bracket can be a single, monolithic piece. While the insert molding of a plastic bracket 60 around a metal striker 30 is a preferred method, other materials may be used and other fabrication and assembly methods may be employed instead, if so desired. For example, the bracket and striker may both be made of metal and/or the bracket formed in multiple pieces that are assembled to the striker. The installation and adjustment of the adjustable striker assembly 20 will now be described. The mounting flanges 62 are placed in the desired location of the door or compartment, as the case may be, and the fasteners (not shown in this embodiment) are installed through the mounting holes 64 to secure the striker assembly 20 in place. If the D-shaped hoop 36 is in the proper location for engagement with the latch, then no further adjustment needs to be made. If the D-shaped hoop 36 does not engage properly with the latch, then the hoop 36 is grasped and rotated one-half turn relative to the bracket 60 —clockwise or counterclockwise depending upon the direction of misalignment. If additional adjustment is needed, then it can be rotated another half turn. When the D-shaped hoop 36 is grasped and a significant level of torque is applied, the torque causes the barb 78 of the flexible arm 76 to be flexed out of the V-groove 48 . This then releases the main shaft 32 to continue rotation through one hundred eighty degrees. At that time, the barb 78 aligns with the opposite V-groove 48 , allowing the arm 76 to flex back and bias the barb 78 into this opposite V-groove 48 . This will hold the striker 30 in its new position and correct orientation relative to the bracket 60 . One will note that the adjustment has been accomplished without the need for any tools. Thus, if misalignment should occur due to installation of an instrument panel (not shown) into a vehicle body (not shown), adjustment can be made quickly and easily. One will also note that the amount of torque required to cause rotation of the striker 30 relative to the bracket 60 may be easily tuned during design of the bracket 60 for a particular application by making small changes to the thickness of the flexible arm 76 and/or changes to the barb 78 and V-groove 48 . FIG. 5 illustrates a second embodiment. This embodiment has many items in common with that of the first embodiment, and to avoid unnecessary repetition of the description, the same reference numerals have been used but falling within the 100-series. The striker assembly 120 is shown mounted to a door or compartment 124 , as the case may be, and held in place with fasteners 126 (only two shown) that engage mounting holes 164 in the mounting flanges 162 . The striker 130 may be the same as in the first embodiment, if so desired, while the bracket 160 has changed somewhat. In this embodiment, the side support walls and rear support walls have been eliminated. So, the trunnion plate 168 , end wall 170 , and flexible arm 176 extend directly from the mounting flanges 162 . Also, the edges of the mounting flanges 162 now define the central opening 166 . The fabrication of the striker assembly 120 and method of adjustment may be the same as in the first embodiment. FIGS. 6-9 illustrate a third embodiment. This embodiment has many items in common with that of the first embodiment, and to avoid unnecessary repetition of the description, the same reference numerals have been used but falling within the 200-series. For this striker assembly 220 , the striker 230 may again be the same as in the first embodiment, if so desired, while the bracket 260 has changed. The flexible arm 276 , mounting flanges 262 , side support walls 265 , trunnion plate 268 and rear support walls 274 may be the same as in the first embodiment. In this embodiment, however, the end wall 270 is canted to match the angle of the threads 242 on the main shaft 232 . While this end wall 270 and threaded hole 272 may possibly be easier to mold around the threads 242 than the first embodiment, it may induce more friction between the threads 242 and the threaded hole 272 than the configuration in the first embodiment. While certain embodiments of the present invention have been described in detail, those familiar with the art to which this invention relates will recognize various alternative designs and embodiments for practicing the invention as defined by the following claims.

A striker assembly, and method of forming, for use with a closure for a vehicle compartment is disclosed. The striker may have a main shaft with a hoop extending from a first end, threads extending around the main shaft, and a rotational retention feature adjacent to the threads. A bracket may have a pair of mounting flanges, a first wall including a threaded hole therethrough that engages the threads, a second wall, spaced from the first wall, having a hole therethrough that engages the main shaft, and a flexible arm having a catch engageable with the rotational retention feature.

BACKGROUND OF THE INVENTION 1. Field of the Invention The invention is related generally to the field of acquisition and interpretation of measurements made by well logging instruments for the purpose of determining the properties of earth formations. More specifically, the invention is related to a method for making resistivity measurements to measure the effects of invasion of borehole fluid into the formation. 2. Background of the Art The estimation of hydrocarbon reserves depends heavily on the accuracy of resistivity data and the reliability of their interpretation. One of the primary difficulties in formation analysis from borehole surveys is the need to determine and compensate for the effects of invasion. Invasion takes place in porous permeable zones where the hydrostatic/dynamic pressure of the drilling mud is greater than the formation pore pressure. The invasion of the mud filtrate will cause a radial variation of the formation resistivity. One of the objectives of resistivity measurements is to get an estimate of the resistivity of the uninvaded formation. Such a resistivity estimate can be used as a basis for determining hydrocarbon saturation in the formation, and thus serve as a basis for estimating total recoverable hydrocarbons in place. See, for example, U.S. Pat. No. 5,883,515 to Strack et al., having the same assignee as the present invention and the contents of which are incorporated herein by reference. In addition to estimating recoverable hydrocarbons, another parameter of interest is the formation permeability. The formation permeability is related to the rate at which reservoir is capable of producing, a very important factor in determining potential profitability and cash flow from a reservoir. Reservoir permeability can be estimated using different techniques, such as core analysis, well-log correlation, and well testing. These different techniques result in the values of permeability representative of different sample volumes of the reservoir. Core analysis is expensive and time-consuming and the sample size is too small to characterize a reservoir. Well testing is also a time consuming procedure. Another method that may be used is to make time-lapse measurements of resistivity. The rate at which the borehole fluid invades the formation depends on the permeability; consequently, the radial variation of resistivity at different times will depend on the permeability. Simulated resistivity logs based on petrophysical models may be compared to actual resistivity profiles, and model parameters such as permeability may be adjusted to provide a good match between the simulated and actual measurements. See, for example, U.S. Pat. No. 5,497,321 to Ramakrishnan et al. and Frenkel (SPE 100359). Having discussed some examples on the utility of making resistivity measurements, we briefly discuss some prior art methods of determining formation resistivity. Frenkel et al. (SPE 62912) discusses the use of galvanic measurements (specifically, High-Definition Lateral Log measurements) for determining a resistivity model. FIG. 3 illustrates resistivity model used in Frenkel. The model includes formation layers surrounding a borehole having a diameter D bh filled with mud having resistivity R m . This model is typical of those that have been used in the past, and basically approximates each layer by an invaded zone having a length L x0 and resistivity R x0 and an uninvaded zone having a resistivity R t . A possible refinement that has been used is to add a flushed zone, defining a three radial-layer model. A review of Frenkel shows that a typical depth of the invaded zone is about 2 ft. (0.6 m). Similar models have been used for inverting resistivity measurements made with an induction logging tool. See, for example, Semmelbeck et al. (U.S. Pat. No. 5,663,499). The accuracy of saturation estimates and permeability estimates depends on the accuracy and reliability of methods of obtaining and analyzing formation resistivity measurements to provide an accurate radial resistivity model. In this regard, it is desirable to be able to define multiple zones with high resolution in the first 2 ft. (0.6 m) or so. The present invention satisfies this need. SUMMARY OF THE INVENTION One embodiment is a method of evaluating invasion of a borehole fluid into an earth formation. The method includes conveying into a borehole a resistivity measuring instrument having at least one transmitter and a plurality of pairs of receivers spaced apart from the at least one transmitter. The transmitter is activated and at least one frequency and signals are induced in the plurality of pairs of receivers that are indicative of the resistivity of at least three subzones inside an invaded zone. The induced signals are filtered to provide an estimate of the resistivity of the at least three subzones. The estimate of the resistivity may be recorded on a suitable medium. One of each of the pairs of receivers may be used as a main coil and the other of each of the pairs of receivers may be used as a bucking coil. The at least one frequency may be less than 25 kHz. The method may further include activating the at least one transmitter at a plurality of frequencies greater than about 0.5 MHz and inducing additional signals in the receivers, and using the estimate of the resistivity and the additional signals for providing an updated estimate of the resistivity of the at least three zones. The borehole fluid may have a resistivity greater than about 0.10 Ω-m. The at least three zones may include at least five zones. A distance from the transmitter to at least one of the receivers may be greater than 1 m. A distance from a transmitter to at least one of the receivers may be less than 0.3 m. Providing the updated estimate of resistivity may be done using an iterative gradient method. The measurements may be repeated after an elapsed time interval, and the resistivities before and after the elapsed time interval may be used for estimating a permeability of the formation. Another embodiment of the invention is an apparatus for evaluating invasion of a borehole fluid into an earth formation. The apparatus includes a resistivity measuring instrument configured to be conveyed into the borehole. The instrument includes at least one transmitter configured to be activated at least one frequency, and a plurality of pairs of receivers spaced apart from the at least one transmitter. The receivers are configured to receive signals resulting from activation of the at least one transmitter. The signals are indicative of a resistivity of at least three subzones inside an invaded zone. The apparatus further includes a processor configured to filter the signals from the plurality of pairs of receivers and provide an estimate of the resistivity of the at least three zones. The processor may be further configured to record the estimate of the resistivity on a suitable medium. Each of the pairs of receivers may include a main receiver and a bucking receiver. The at least one frequency may be less than 25 kHz. The at least one transmitter may be further configured to be activated at a plurality of frequencies greater than about 0.5 MHz, and the processor may be further configured to use the estimate of the resistivity and additional signals from the plurality of receivers resulting from the further activation of the at least one transmitter to provide an updated estimate of the resistivity of the at least three zones. The borehole may contain a fluid having a resistivity greater than about 0.1 Ω-m. The at least three zones may include five zones. A distance from the transmitter to at least one of the pairs of receivers may be greater than 1 m. A distance from the transmitter to at least one of the pairs of receivers may be less than 0.3 m. The processor may be configured to provide the updated estimate of the resistivity using an iterative gradient method. The apparatus may include a conveyance device a configured to convey resistivity measuring instrument into the borehole; the conveyance device may be a wireline or a drilling tubular. Another embodiment of the invention is a computer-readable medium for use with an apparatus for evaluating invasion of a borehole fluid into an earth formation. The apparatus includes a resistivity measuring instrument configured to be conveyed into the borehole. The instrument includes at least one transmitter configured to be activated at least one frequency and a plurality of receivers spaced apart from the at least one transmitter configured to receive signals resulting from the activation of the at least one transmitter. The signals are indicative of a resistivity of at least three subzones inside an invaded zone. The medium includes instructions which enable a processor to filter the signals from the plurality of receivers and provide an estimate of the resistivity of the at least three subzones. The medium may include a read-only memory, a programmable read-only memory, an electrically programmable read-only memory, an electrically alterable read-only memory, an electrically erasable and programmable read-only memory, a flash memory, an optical disk, a hard drive, an iPod®, and/or a non-volatile read-write memory. BRIEF DESCRIPTION OF THE FIGURES The present invention is best understood with reference to the accompanying figures in which like numerals referred to like elements and in which: FIG. 1 shows an induction instrument disposed in a wellbore penetrating an earth formation; FIG. 2 shows the arrangement of transmitter and receiver coils in an embodiment of the present disclosure, FIG. 3 shows an exemplary model with an invaded zone; FIGS. 4A , 4 B show an exemplary model with biopolymer drilling mud, water-bearing formation. R m =0.1 Ω-m; R t =4 Ω-m. FIG. 4 a : invasion zone is represented by 3 radial zones, FIG. 4B : invasion zone is represented by 5 radial zones; FIGS. 5A , 5 B show an exemplary model with water based mud, water-oil-bearing formation. R m =2 Ω-m; R t =10 Ω-m. In FIG. 5A , the invasion zone is represented by 3 radial zones and in FIG. 5B , the invasion zone is represented by 5 radial zones; FIGS. 6A , 6 B show an exemplary model with Oil-based mud, oil-bearing formation. R m =1000 Ω-m; R t =20 Ω-m. In FIG. 6A , the invasion zone is represented by 3 radial zones and in FIG. 6B , the invasion zone is represented by 5 radial zones; FIGS. 7A , 7 B show inversion results for the model of FIG. 4B using low frequencies ( FIG. 7A ) and high frequencies ( FIG. 7B ); FIGS. 8A , 8 B show inversion results for the model of FIG. 5A using low frequencies ( FIG. 8A ) and high frequencies ( FIG. 8B ); FIGS. 9A , 9 B show inversion results for the model of FIG. 5A using low frequencies ( FIG. 9A ) and high frequencies ( FIG. 9B ); FIG. 10 shows an inversion result for the model of FIG. 5B ; FIG. 11 shows an example of geometric factors for three layers in the invaded zone; and FIG. 12 shows filters for signal transformation, when the signal is sensitive to only one of the three layers in turn. DETAILED DESCRIPTION OF THE INVENTION Referring now to FIG. 1 , an induction logging tool 20 suitable for use with the present invention is shown positioned in a borehole 22 penetrating earth formations 54 . The tool 20 , which is suspended in the borehole 22 by means of a wireline cable 24 , includes a borehole sonde 34 and an electronic circuitry section 32 . The tool 20 is lowered into the borehole 22 by a cable 24 , which passes over a sheave 31 located at the surface of the borehole 22 . The cable 24 is typically spooled onto a drum 30 . The cable 24 includes insulated electric conductors for transmitting electrical signals. The electronic circuitry section 32 of the tool 20 receives signals from the sonde 34 to perform various analog and digital functions, as will be described later. The sonde 34 may include a plurality of coils 40 - 52 . Coil 46 is a transmitter coil for transmitting an oscillating signal into the adjacent surrounding geological formation 54 . It is contemplated that any of a number of oscillating voltage signals having multiple frequency components can be used. Further, it is desirable that, on occasion, a single-frequency signal, such as a sinusoidal signal, is used. The oscillating voltage signal applied to the coil 46 generates a current in coil 46 which in turn generates an electromagnetic field in the surrounding formation 54 . The electromagnetic field, in turn, induces eddy currents, which flow coaxially with respect to the borehole 22 . The magnitudes of the eddy currents are related to the conductivity of the surrounding formation 54 . The remaining coils 40 , 42 , 44 , 47 , 48 , 50 and 52 are receiver coils in which signals are induced by the electric fields caused by the eddy currents produced in the formation. As the tool 20 is raised in the borehole 22 , the conductivity of the surrounding formation 54 can be determined from the received signals in order that a bed or layer 55 having a conductivity that is indicative of the possibility of containing hydrocarbons may be located. The number of receiver coils shown is for illustrative purposes only; more or fewer coils may be used. The electronic circuitry section 32 may include a converter circuit 60 , a stacker circuit 62 , a random access memory (RAM) 63 , and a telemetry circuit 61 . The converter circuit 60 comprises a plurality of pre-amplifiers, filters, and analog-to-digital (A/D) converters for receiving signals from the receiver coils 40 - 52 and transforming them into digitized signals for further processing by the stacker circuit 62 . The analog voltage signals provided by the receiver coils 40 - 52 are digitally sampled according to a predetermined sampling rate in the period defined by the fundamental frequency of the transmitter signal, which in a typical embodiment is approximately 10 kHz. The sampling may be repeated over a large number of transmitter voltage signal cycles, preferably at least 1,024 cycles to improve the signal-to-noise ratio of the received signals. To reduce the amount of data that must be stored or transmitted, corresponding digital samples taken in each of the transmitter cycles are summed. The summed digital signal samples corresponding to each of the plurality of receiver coils form corresponding stacked signal samples, which are stored in the RAM 63 . The stacked signals corresponding to the plurality of receiver coils 40 - 52 can then be retrieved from the RAM 63 and can be transmitted by the telemetry circuit 61 through the cable 24 to a processor 64 which forms part of the surface equipment 26 , where analyses of the stacked signals can be performed. Alternatively, processing of at least part of the data could be performed downhole using a processor at a suitable location (not shown) and results of the processing telemetered uphole. In an alternative embodiment, a processor having sufficient digital signal processing capabilities could form part of the electronic circuitry section 32 . Thus, it is contemplated that the required discrete Fourier transform could be performed downhole, which would further reduce the amount of data to be transmitted to the surface. Turning now to FIG. 2 , a suitable configuration of transmitter and receiver coils is illustrated. Three-coil induction tool configuration is considered. Its elements are: transmitter 101 of moment M z , main receiver 103 (Z 1 ) and bucking coil 105 (Z 2 ). Thus, each of the receiver coils of FIG. 1 consists of the main receiver and a bucking coil. The tool is positioned on the axis of the borehole 107 of resistivity R m and radius r m . Resistivity of the formation is R t . Radial resistivity profile in the invasion zone 109 of radius r xo is described by a stepwise function R xo (r). The ability of the methods of the present invention to respond to a variety of different models was evaluated. These were intended to cover a variety of drilling muds and formation resistivities. An exemplary model, shown in FIGS. 4 a , 4 b simulated a biopolymer drilling mud, water-bearing formation. R m =0.1 Ω-m; R t =4 Ω-m. In FIG. 4 a , the invasion zone is represented by 3 radial zones while in FIG. 4 b , the invasion zone is represented by 5 radial zones. The model of FIGS. 5 a , 5 b are for water based mud, water-oil-bearing formation. R m =2 Ω-m; R t =10 Ω-m. In FIG. 5 a , the invasion zone is represented by 3 radial zones while in FIG. 5 b , the invasion zone is represented by 5 radial zones. FIGS. 6 a , 6 b show an exemplary model with Oil-based mud, oil-bearing formation. R m =1000 Ω-m; R t =20 Ω-m. In FIG. 6 a , the invasion zone is represented by 3 radial zones and in FIG. 6 b the invasion zone is represented by 5 radial zones. Results from three types of data acquisition and inversion are discussed. The first, referred to as the Low Frequency Method, used a dense set of arrays spacings and a single frequency. The arrays were chosen with z 2 =0.25, 0.562, 0.875, 1.188, 1.50, 1.812, 2.125 (m), z 1 =0.8z 2 ; f=10 kHz; M=7−the number of measurements. The measured signal is the difference of the emf real parts in the three-coil array. The second method discussed, referred to as the High Frequency Method, used frequencies up to 10 MHz. This is an attempt to avoid the dielectric effect and technological complications of constructing coils at very high frequencies. The same array spacings as in the Low Frequency method were used, and the following set of frequencies were selected: f=2.128, 3.275, 4.308, 5.237, 6.073, 6.826, 7.503, 8.113, 8.662, 9.156, 9.6, 10 (MHz). This gives a total of 84 measurements. The third method discussed is the Very High Frequency Method. The following system of lengths and frequencies was chosen: z 2 =0.25, 0.494, 0.732, 0.964, 1.19, 1.41, 1.624, 1.832 (m); z 1 =0.2, 0.4, 0.6, 0.8, 1.0, 1.2, 1.4, 1.6 (m); f=20.988, 50.791, 74.633, 93.706, 108.96, 121.17, 130.94, 138.75, 145, 150 (MHz); M=80. The measured signal is the phase difference in the three-coil array. An assumption was made that the dielectric constant does not have a significant influence on the resolution, and it has been assumed that it is equal to zero. For each of the methods, the results of inversion are shown. Simulated data {right arrow over (F)} E are modeled based on the true model, with an additive random error (less than or equal to 1% at each measurement). The pattern of error distribution is unknown. We also assume that the resistivities of the mud and formation are known. The mud resistivity may be estimated, for example, using the apparatus described in U.S. Pat. No. 6,801,039 to Fabris et al. the contents of which are incorporated herein by reference. The formation resistivity may be obtained using low frequency measurements made with an induction tool having a large transmitter-receiver distance. To estimate the quality of the solution we will use mean square deviation of model data {right arrow over (F)} (based on the resulting model) from synthetic data: Δ = 1 M ⁢ ( F j E - F j F j E ) 2 · 100 ⁢ % , and also we will calculate the integral conductivity of invaded zone. Turning now to FIG. 7 a , results for the model of FIG. 4 a are shown, Shown is the actual model 700 along with an inverted model 702 derived using the Low Frequency Method. For this example, the value of Δ=0.5%. The relative error of defining the first subzone resistivity is 3.2%, of the second zone—7.6%, of the third zone—15%. The true integral conductivity is 0.567 S/m, while the recovered conductivity is 0.561 S/m, giving a relative error of 1%. The results of the low frequency method were used as an initial estimate in the high frequency method. FIG. 7 b shows the result of using the High Frequency Method. This has a value of Δ=1%. The relative error of defining the first and the second subzones resistivity is less than 1%, of the third—2.8%. The true integral conductivity is 0.567 S/m, recovered conductivity is 0.570 S/m, with a relative error is 0.5%. Turning now to FIG. 8A , results for the model of FIG. 5A are shown, Shown is the actual model 800 along with an inverted model 802 derived using the Low Frequency Method. For this example, the value of Δ=0.4%. Relative error of defining the first subzone resistivity is 10%, of the second—23%, of the third—2%. The true integral conductivity is 0.108 S/m, recovered conductivity is the same. The results of the low frequency method were used as an initial estimate in the high frequency method. FIG. 8B shows the result of using the High Frequency Method. This has a value of Δ=0.9%. Relative error of defining the first subzone resistivity is 1%, of the second zone —5%, and of the third zone—0.5%. The true integral conductivity is 0.108 S/m, recovered conductivity is the same. Turning now to FIG. 9 a , results for the model of FIG. 6 a are shown, Shown is the actual model 900 along with an inverted model 902 derived using the Low Frequency Method. For this example, the value of Δ=0.4%. Relative error of defining the first subzone resistivity is 27%, of the second zone—16%, and of the third zone—3.3%. The true integral conductivity is 0.016 S/m, and the recovered conductivity is the same. The results of the low frequency method were used as an initial estimate in the high frequency method. FIG. 9 b shows the result of using the High Frequency Method. This has a value of Δ=1%. Relative error of defining the first subzone resistivity is 0.1%, of the second—5%, of the third—2%. The true integral conductivity is 0.016 S/m, recovered conductivity is the same. The following points may be noted with respect to FIGS. 7 a , 7 b , 8 a , 8 b , 9 a and 9 b . The low frequency method gives a higher error with resistive formations and/or resistive mud, while the high frequency method does not appear to have this problem. Compare FIG. 8 a to 8 b , and FIG. 9 a to 9 b . Both methods give an integral conductivity that is close to the true value. The results also show that using the high frequency method, it is possible to resolve three invaded zones within a distance of 0.6 m from the borehole for a wide range of mud, formation and invaded zone properties. Turning next to FIG. 10 , results of using the very high frequency method for the model of FIG. 5 b are shown. The curve 902 is the initial model and 900 is the actual model that is indistinguishable from the final result. Δ=0.5%. The relative error of defining all the subzones resistivity is less than 1%. The true integral conductivity is 0.09 S/m, and the recovered conductivity is the same. Next, we explain in some detail the method of analysis and inversion of the measurements. We consider different linear representations of the forward problem and filters as the mean of solving the inverse problem. Then we use two approaches to determine the resistivities of the invaded zone. The first approach is based on Doll approximation and the second is based on the Jacobian. The first approach allows us to build an initial guess. It is a rigorous solution of the linear system and using it we can distinguish no more than 3 subzones when the resistivities of the borehole and formation are known. The second approach enables to improve the initial guess at high frequencies (less than 10 MHz) and distinguish up to 5 subzones at very high frequencies (less than 150 MHz). It is an iterative procedure which requires a good initial guess. The forward problem of electromagnetic logging can be expressed as: {right arrow over (F)}≈{right arrow over (F)} 0 +Ĝ ( {right arrow over (p)}−{right arrow over (p)} 0 )  (1) or {right arrow over (F)}≈Ĝ{right arrow over (p)}  (2) where G ji is the contribution of parameter p i into the signal F j , i=1, . . . , N, j=1, . . . , M; N is the number of parameters, M is the number of measurements. In this document we use the following linear representations of the forward problem: (a) Doll approximation: G ji is a geometric factor (a value that does not depend on medium conductivity), that is, a contribution into the signal F j from the area corresponding to conductivity σ i when the rest of the medium has zero conductivity, and the area's conductivity is equal to σ 0i , p i =σ i /σ 0i ; (b) the first term in Taylor series: G ji is a derivative of the signal F j with respect to the parameter of the model p i at {right arrow over (p)} 0 . To solve the inverse problem, filters may be used to minimize background parameters effect and to maximize the influence of parameters to be determined. We partition the unknown parameters into N1 background parameters and N2 parameters to be determined. Thus, p i , i=1, . . . , N1 be background parameters; p i , i=N1+1, . . . , N being parameters to be determined, N2=N−N1. We should find the vector {right arrow over (c)} (filter) which satisfies the equation: { g -> l T · c -> = 0 , l = 1 , N ⁢ ⁢ 1 g -> k T · c -> = α k , k = 1 , N ⁢ ⁢ 2 ( 3 ) Here {right arrow over (g)} l —columns of matrix Ĝ corresponding to background parameters, {right arrow over (g)} k —columns corresponding to parameters we intend to determine. We can express the same in a matrix form: Ĝ T {right arrow over (c)}={right arrow over (A)}  (4) The first N1 elements of the vector {right arrow over (A)} are zeros, the next N2 elements are coefficients α k . In this case signal transformation {right arrow over (c)} T {right arrow over (F)} will have reduced sensitivity to parameters p l , l=1, . . . , N1 and increased sensitivity to parameters p k , k=1, . . . , N2 depending on the ratio of α k values. We solve the inverse problem using this transformation. We assume that M≧N and solve the set of equations (4) by means of singular value decomposition (SVD) method. In the case of no singularity at R=N (where R is the rank of the Ĝ matrix) the solution {right arrow over (c)} has the minimal norm and belongs to the subspace {right arrow over (g)} k , k=1, . . . , N2. In the case of an incompatible set, we arrive at an approximate solution at R<N, and this solution corresponds to both the norm minimum and the minimum of quadratic residual. In an alternate approach we solve the inverse problem with Ĝ being the matrix of geometric factors We write equation (2) in normalized form: {right arrow over (F)} N ≈Ĝ{right arrow over (p)}  (2a) where F Nj = F j ⁡ ( σ -> ) F j ⁡ ( σ -> 0 ) , ⁢ G ji = F j ( 0 , … ⁢ , 0 , σ 0 ⁢ i i , 0 , … ⁢ , 0 ) F j ⁡ ( σ -> 0 ) , ⁢ p i = σ i σ 0 ⁢ i . Let the system of measurement consist of three-coil arrays lined up along the borehole axis. The source is a vertical magnetic dipole M z , current's density in which follows the law of e −iωt (f=ω/(2π) is the frequency); the measured signal is the difference of imaginary components of the magnetic field H z in two-coil arrays. Mf of frequencies and Ml of array lengths are used in such a way that Mf·Ml=M. The model of the medium is cylindrically layered. The first layer is the borehole itself with its conductivity σ 1 and radius r 1 , the last layer is the formation of conductivity σ N . Layers numbered from 2 to N−1 constitute invaded zone. Conductivities of layers are σ i , their radii are r i , i=2, N−1; r 1 =0.1 m is the borehole radius; r N−1 =0.6 m is the radius of invaded zone. To analyze how much layers (subzones) the invaded zone can be subdivided into, we assume that resistivities of the borehole and of the formation are known. If this is not so, we should first determine their values. From now on, when we talk about the matrix Ĝ, we will mean only that part of it which corresponds to resistivities of the invaded zone. Let us take the following frequency and lengths: 10 (KHz); 0.20, 0.25, 0.45, 0.562, 0.70, 0.875, 0.95, 1.188, 1.20, 1.50, 1.45, 1.812, 1.70, 2.125 (m), M=Ml=7. As can be seen, the ratios of the first distance to the second distance (and the third to the fourth, fifth to the sixth etc.) is 0.8. There are seven measurements in all when three coil arrays are used. For a three-coil array two consecutive lengths are taken, and the moments are chosen in such a way as to compensate the direct field. The lengths ratio for the three-coil array is 0.8, the length of the next array is 0.25 m greater (according to the short array). For these values and exemplary models of the formation, the relative error for the forward problem equation (2) does not exceed 0.1% (this error increases with frequency and array length). Geometric factors of invaded zone and the borehole are linearly dependent on frequency; and when invaded zone parameters are being determined at a given resistivity of the formation, it is enough to take one frequency and (to reduce equivalence) many lengths. It is known that when the inverse problem is being solved for representation eqn. (2) (or for (2)), if equal relative measurement errors ε and equal solution errors δ are allowed) the relative solution error is defined by this inequality: δ =  p -> Δ   p ->  ≤ cond ⁢ ⁢ G ^ ⁢  F -> Δ   F ->  = cond ⁢ ⁢ G ^ · ɛ ( 5 ) Condition cond Ĝ is defined as following: cond ⁢ ⁢ G ^ = λ max λ min , where λ max =λ N and λ min are the maximal and the minimal (non-zero) eigenvalues of the information matrix Â=Ĝ T Ĝ; ∥{right arrow over (p)} Δ ∥, ∥{right arrow over (F)} Δ ∥ are absolute errors. Experience shows that for satisfactory inversion it is necessary for the matrix Ĝ to contain a complete rank (λ min =λ 1 ) and cond Ĝ≦1/ε. In a general case the number of measurements may be reduced, the major reduction criterion being the decrease of the value of cond Ĝ due to removing linearly dependent matrix Ĝ rows. By means of this reduction it is possible to greatly improve the estimated error eqn. (5). Another way of improving the ones, that is, of decreasing cond Ĝ, is to choose adequate normalization, which is achieved by multiplying the matrix of the forward problem by the appropriate diagonal matrix from the left. Results of constructing the matrix Ĝ have shown that for subdivision of the invaded zone into layers of equal thickness we obtain the following: cond Ĝ=1525 for five layers, cond Ĝ=117 for four layers, cond Ĝ=14.5 for three layers. Thus, for a 1% measurement error, invaded zone may be subdivided into no more than three subzones, and the quality of the obtained solution will depend on the model. For a model in which r 1 =0.1 m, r 2 =0.266 m, r 3 =0.433 m, r 4 =0.6 m; {right arrow over (σ)} 0 =(1, . . . 1), the layers are shown, for example, in FIGS. 4A , 4 B, 5 A, 5 B, 6 A and 6 B. Geometric factors for the three uniform layers inside invaded zone are shown in FIG. 11 as 1001 , 1003 and 1005 . It should be noted that the borehole's geometric factor is close to zero (less than 0.04), while the formation's geometric factor ranges from 0.3 to 1, therefore we need to know an accurate value of the formation's conductivity, while the value of the borehole conductivity may be approximate. Filters for signal transformation, when the signal is sensitive to each of the three layers in turn are shown in FIG. 12 by 1201 , 1203 , 1205 . Filters may be calculated from this formula: c -> k = ∑ i = 1 N 2 ⁢ v ki λ i ⁢ u -> i where v ki is k-component of the eigenvector of the matrix Â=Ĝ T Ĝ, corresponding to the eigenvalue λ i ; {right arrow over (u)} i is the eigenvector of the matrix  T =ĜĜ T , corresponding to the same eigenvalue, N 2 =N−2. When the conductivities of the borehole and the formation are known (R 0 and R t ) the forward problem of eqn. (2) may be written out like this: {right arrow over (F)} N −{right arrow over (F)} N0 ≈Ĝ{right arrow over (p)}   (2b) where {right arrow over (F)} N0 ={right arrow over (F)} N (σ 1 , 0, . . . , 0)+{right arrow over (F)} N (0, . . . 0, σ N ), p i = σ i σ 0 ⁢ i , ⁢ i = 2 , N - 1. Let {right arrow over (c)} k be a filter constructed for defining the k-th layer of invaded zone. Then c → k T ( F → N - F → N ⁢ ⁢ 0 ) ≈ σ k + 1 σ 0 ⁢ k + 1 . Substituting synthetic data {right arrow over (F)} E for {right arrow over (F)}, we obtain the value of σ k+1 . For eqn. (2b) the estimated solution error can be presented as follows: δ ≤ cond ⁢ G ^ ·  F → N ⁢ ⁢ Δ   F → N - F → N ⁢ ⁢ 0  = cond ⁢ G ^ · ɛ ·  F → N ⁢   F → N - F → N ⁢ ⁢ 0  ≤ cond ⁢ G ^ · ɛ · 1 1 -  F → N ⁢ ⁢ 0  /  F → N  ( 5 ⁢ a ) That is, the normalized signal decrease also influences the error of the solution. Synthetic data {right arrow over (F)} E can be modeled based on the true model, and a random error added (less than or equal to 1% at each measurement). The pattern of error distribution is unknown. To estimate the quality of the solution we will use mean square deviation of model data from synthetic data (in %) Δ = 1 M ⁢ ( F j E - F j F j E ) 2 · 100 ⁢ ⁢ % and the mean difference of filtered values Δ k = c → k T ⁢ F → N E - c → k T ⁢ F → N c → k T ⁢ F → N E · 100 ⁢ ⁢ % , ⁢ F Nj E = F j E ⁡ ( σ → ) F j ⁡ ( σ → 0 ) . The matrix of the system does not depend on medium conductivities so the inverse problem solved by filters is equivalent to solving a set of equations by SVD method. One might draw certain conclusions here. It appears impossible to evaluate high values of resistivities due to their small contribution into the signal. However, this result is helpful, being an initial model for further detailed interpretation, e.g. using filters built for higher frequencies and depending on the model itself. An important feature of the inversion method discussed here is instantaneous computing of conductivities based on given values of {right arrow over (c)} k . The values of {right arrow over (c)} k T {right arrow over (F)} N0 can be computed by means of linear interpolation across a certain grid {σ 1 ,σ N }. The quality of inversion depends on the value of 6 in the estimation given by eqn. (5a). Thus, when δ gets too high (over 100%), one could try reducing the dimension of the set or adding new measurements to reduce equivalence. Generally equivalence is defined by the inequality ∥d{right arrow over (F)}∥≦ε, where dF j = F ~ j - F j F j ; F j - data , corresponding to inverted model, {tilde over (F)} j −data, corresponding to equivalent model. An important advantage of the solution obtained is that now one can easily compute the integral conductivity of invaded zone. We next discuss solving inverse problem with Ĝ being a Jacobian. Eqn. (1) may be rewritten as follows: {right arrow over (F)}≈{right arrow over (F)} 0 +Ĝ{right arrow over (p)}   (1a) Elements of the matrix Ĝ are derivatives of the signal with respect to the logarithm of conductivity computed for the initial model with the conductivities {right arrow over (σ)} 0 : G ji = - ∂ F j ∂ σ i + 1 ⁢ ( σ → 0 ) · σ 0 ⁢ ⁢ i + 1 = ∂ F j ∂ ρ i + 1 ⁢ ( ρ → 0 ) · ρ 0 ⁢ i + 1 , ⁢ p i = ρ i + 1 - ρ 0 ⁢ i + 1 ρ 0 ⁢ i + 1 ( 6 ) The arrays lengths are the same as for geometric factors. Frequencies are limited from above by 10 MHz (this is an attempt to avoid the permittivity effect and hardware complications related to constructing coils at very high frequencies). Let us try to find out how many subzones could be identified within the invaded zone in this case, and evaluate the quality of the solution. For the eqn. (1) in light of eqn. (6), the solution error can be defined as follows:  δ →  ≤ ɛ λ 1 ⁢  F → 0  or, if the solution contains identical relative errors, δ ≤ ɛ λ 1 ⁢ N ⁢  F → 0  . ( 7 ) The value of cond Ĝ shows how admissible domains of changing the generalized parameters differ. Because at high frequencies signals go through zero, derivatives of the logarithm of the signal are not used. To improve the estimation (7) and to reduce the value of cond Ĝ, normalization to the norm of the corresponding row of the matrix Ĝ may be useful: F → N - F → N ⁢ ⁢ 0 ≈ G ^ N ⁢ p → , ⁢ F Nj = F j ⁡ ( σ → ) g j , ⁢ g j = ∑ i = 1 N 2 ⁢ ⁢ G ji 2 , ⁢ G Nji = G ji g j . From now on, we will drop the subscript N for Ĝ N : {right arrow over (F)} N −{right arrow over (F)} N0 ≈Ĝ{right arrow over (p)}   (1-b) To estimate the error for identical relative errors, the following formula is used: δ ≤ ɛ λ 1 ⁢ N ⁢  F → N ⁢ 0  ( 7 ⁢ - ⁢ a ) Computations for different models have shown that, in this system of measurement, only high frequencies (over 1 MHz) are important. Optimization criteria were the minimal values of cond Ĝ and δ from eqn. (7a). According to these criteria for different model types the following system of frequencies and lengths was chosen: 2128.1, 3275.3, 4307.8, 5237.0, 6073.3, 6826.0, 7503.4, 8113.0, 8661.7, 9155.6, 9600.0, 10000.0 (KHz); 0.20, 0.25, 0.45, 0.562, 0.70, 0.875, 0.95, 1.188, 1.20, 1.50, 1.45, 1.812, 1.70, 2.125 (m), for a total of M=84 measurements. It should be noted that this choice of lengths and frequencies is not the best for all models. It appears reasonable to have a finer grid (system of measurement) for both lengths and frequencies, so that for each model type the optimal values could be chosen. In Table A-1 values of cond Ĝ and δ for the exemplary models are given. The invaded zone is subdivided into three or four equal subzones. The solution of eqn. (7a) is not as critical as the solution of eqn. (5a) for the exact problem because the forward problem is not linear, and inversion is an iterative process. Whether or not it is possible to identify three or four layers could be determined, for some model classes, only after inversion tests. TABLE A-1 ρ 1 , Three layers Four layers Ohm · m cond Ĝ Δ cond Ĝ δ 0.1 6.9 17 52 115 2.0 10 26 75 162 1000.0 26 72 52 231 In the case when initial approximation for invaded zone's conductivity is not available or is given with a large error, the error of the linear representation of forward problem of eqn. (1b) is also large. For this reason, an iterative process is used for inversion. First of all the sequence of determining parameters is fixed. For example first determining parameter—those which corresponds to maximal diagonal element of matrix Â, last determining parameter corresponds to minimal diagonal element of Â. To evaluate each parameter, a filter {right arrow over (c)} k is created, and the value Δ k is minimized on the interval of valid values of the parameter. The conductivities {right arrow over (σ)} 0 are formed each time as prior values, already known at every previous step of minimization. When the determination of all unknown parameters of the model is complete, there is possible to repeat all the previous steps (to perform the second iteration), based on the model obtained after the first iteration. The second iteration could be necessary when initial approximation is too far from the true value and the Δ≦∥δ{right arrow over (F)} E ∥·100% has not been achieved yet. When minimization on the admissible interval is performed, several minimums are possible. In this case, the minimum with the smallest Δ is chosen. In the case when Δ is small enough, a reverse of Ĝ by means of SVD could be performed, the criterion for success may be the Δ decrease. The condition for ending the inversion process is stabilization of the solution over the subsequent iterations. The results of inversion by geometric factors can be greatly improved at the frequencies up to 10 MHz. In some cases it is possible to distinguish 4 subzones. At the same time initial guess may provide only two subzones. It may occur when formation is much more conductive relatively invaded zone. Then special formulae are applied to construct missing initial resistivities. Formulae are based on the principle of saving integral conductivity of invaded zone. First case is when from two initial resistivities we obtain three ones. Let ρ 01 , ρ 02 be two initial resistivities, ρ 1 , ρ 2 , ρ 3 be three initial resistivities we want to obtain. Then ρ 1 =ρ 01 , ρ 3 = ρ 02 , ⁢ ρ 2 = 2 1 ρ 01 + 1 ρ 0 ⁢ ⁢ 2 Let ρ 01 , ρ 02 be two initial resistivities, ρ 1 , ρ 2 , ρ 3 , ρ 4 be four initial resistivities we want to obtain. Then ρ 1 = ρ 01 , ⁢ ρ 4 = ρ 02 , ⁢ ρ 2 = ρ 3 = 2 1 ρ 01 + 1 ρ 0 ⁢ ⁢ 2 Let ρ 01 , ρ 02 , ρ 03 be three initial resistivities, ρ 1 , ρ 2 , ρ 3 , ρ 4 be four initial resistivities we want to obtain. Then ρ 1 = ρ 01 , ⁢ ρ 4 = ρ 03 , ⁢ ρ 2 = 3 1 ρ 01 + 2 ρ 02 , ⁢ ρ 3 = 3 1 ρ 03 + 2 ρ 02 Computations of different models have shown that, in order to identify five subzones within invaded zone, one needs frequencies of up to 150 MHz. The optimization criteria for the system of observation were the minimal values of cond Ĝ and δ from eqn. (A-7a), just as before. The following system of frequencies and lengths was chosen: 20988.4, 50790.7, 74632.6, 93706.1, 108964.8, 121171.9, 130937.5, 138750.0, 145000.0, 150000.0 (KHz); 0.2, 0.25, 0.4, 0.494, 0.6, 0.732, 0.8, 0.964, 1.0, 1.19, 1.2, 1.41, 1.4, 1.624, 1.6, 1.832 (m), M=80. The measured signal is the phase difference in the three-coil array. Due to the impossibility of considering the permittivity effect at this stage of research, a presupposition has been made that permittivity ε does not have any significant impact on the resolution, and it has been assumed that ε=0. An example of inverting for five subzones is shown in FIG. 10 . For the initial approximation based on a two-layer model of invaded zone the following formulae have been used. ρ 1 = ρ 01 , ⁢ ρ 5 = ρ 02 , ⁢ ρ 3 = 2 1 ρ 01 + 1 ρ 02 , ⁢ ρ 2 = 3 1 ρ 3 + 2 ρ 01 , ⁢ ρ 4 = 3 1 ρ 3 + 2 ρ 02 For initial approximation based on a three-layer model of invaded zone the formulae were: ρ 1 = ρ 01 , ⁢ ρ 3 = ρ 02 , ⁢ ρ 5 = ρ 03 , ⁢ ρ 2 = 3 1 ρ 02 + 2 ρ 01 , ⁢ ρ 4 = 3 1 ρ 02 + 2 ρ 03 For some models and frequencies, the phase difference as function of the length jumps through zero. Such jump functions can cause serious difficulties in inversion, and for this reason the measurement corresponding to a negative phase difference was deleted from the system of measurement. It is necessary to point out the following. A non-conductive borehole has practically no impact on the signal, thus we do not need to know its resistivity precisely. For instance, one can assume it is 900 Ω-m instead of 1000 Ω-m. If the resistivity of the formation is known with a significant error, e.g. 10%, this parameter may be included into the set (A-1) and corrected in the same way as the resistivities of invaded zone. It might be that, to determine six parameters, the system of observation will have to be modified. When the resistivities of the formation and of the conductive borehole are known with poor precision, these parameters could be included into the set (A-1). Then δ would almost double, and cond Ĝ would double or even triple. At the same time, the real area of equivalence would increase greatly, and, to reduce it, the system of measurement will have to be modified. Without this modification, the chosen system of measurement would enable us to determine only six parameters with reasonable reliability. Therefore, preliminary inversion should give a good evaluation of either the borehole resistivity or the resistivity of the formation. Once the formation resistivity has been determined, values of the determined resistivity may be displayed as a log and/or stored on a suitable medium. Those versed in the art would recognize that knowledge of formation resistivity is of great utility in the evaluation and development of hydrocarbon reservoirs. Specifically, accurate resistivity measurements may be used to interpret petrophysical quantities such as water saturation and sand content. Further applications include geologic correlation, determination of hydrocarbon/water contact, invasion profile definition, fracture identification and estimation of moveable hydrocarbons. Thus, with the knowledge of resistivity, decisions can be made about additional evaluation wells to be drilled and location and number of development wells. Time-lapse measurements of resistivity may be used to estimate formation permeability. Implicit in the acquisition and processing the data is the use of a processor. The term processor is intended to include such devices as a field processing gate array (FPGA). The processor may carry out instructions stored on a computer-readable medium such as a read-only memory (ROM), a programmable read-only memory (PROM), an electrically programmable read-only memory (EPROM), an electrically alterable read-only memory (EAROM), an electrically erasable and programmable read-only memory (EEPROM), a flash memory, an optical disk, a hard drive, an iPod®, and/or a non-volatile read-write memory (NOVRAM). While the foregoing disclosure is directed to the preferred embodiments of the invention, various modifications will be apparent to those skilled in the art. It is intended that all variations within the scope and spirit of the appended claims be embraced by the foregoing disclosure.

Measurements made with an induction logging tool are processed to provide a resistivity model of fluid invasion of the formation. Up to five zones can be determined over a radial distance of about 0.6 m. It is emphasized that this abstract is provided to comply with the rules requiring an abstract which will allow a searcher or other reader to quickly ascertain the subject matter of the technical disclosure. It is submitted with the understanding that it will not be used to interpret or limit the scope or meaning of the claims.

CROSS-REFERENCE TO RELATED APPLICATIONS [0001] This application claims the benefit of priority to U.S. provisional application Ser. Nos. 60/718,689, filed Sep. 19, 2005; 60/754,392, filed Dec. 27, 2005; 60/763,593, filed Jan. 30, 2006; 60/752,434, filed Dec. 20, 2005; 60/753,220, filed Dec. 21, 2005; 60/763,696, filed Jan. 30, 2006; and 60/839,947, filed Aug. 23, 2006, herein incorporated by reference. FIELD OF THE INVENTION [0002] The invention encompasses processes for the synthesis of (S)-(+)-3-(aminomethyl)-5-methylhexanoic acid, (S)-Pregabalin, and intermediates of (S)-Pregabalin. BACKGROUND OF THE INVENTION [0003] (S)-Pregabalin, (S)-(+)-3-(aminomethyl)-5-methylhexanoic acid, a compound having the chemical structure, is also known as γ-amino butyric acid or (S)-3-isobutyl GABA. (S)-Pregabalin, marketed under the name LYRICA®, has been found to activate GAD (L-glutamic acid decarboxylase). (S)-Pregabalin has a dose dependent protective effect on-seizure, and is a CNS-active compound. (S)-Pregabalin is useful in anticonvulsant therapy, due to its activation of GAD, promoting the production of GABA, one of the brain's major inhibitory neurotransmitters, which is released at 30 percent of the brains synapses. (S)-Pregabalin has analgesic, anticonvulsant, and anxiolytic activity. [0004] Several processes for the synthesis of (S)-Pregabalin are known. For example, see DRUGS OF THE FUTURE, 24 (8), 862-870 (1999). One such process is illustrated in scheme 1. [0005] In Scheme 1, 3-isobutyl glutaric acid, compound 2, is converted into the corresponding anhydride, compound 3, by treatment with refluxing acetic anhydride. The reaction of the anhydride with NH 4 OH produces the glutaric acid mono-amide, compound 4, which is resolved with (R)-1-phenylethylamine, yielding the (R)-phenylethylamine salt of (R)-3-(carbamoylmethyl)-5-methylhexanoic acid, compound 5. Combining the salt with an acid liberates the R enantiomer, compound 6. Finally, a Hoffmann degradation with Br 2 /NaOH provides (S)-Pregabalin. A disadvantage of this method is that it requires separating the two enantiomers, thereby resulting in the loss of half the product, such that the process cost is high. [0006] Several stereoselective processes for the synthesis of (S)-Pregabalin have been disclosed. For example, U.S. Pat. No. 5,599,973 discloses the preparation of (S)-Pregabalin using stoichiometric (+)-4-methyl-5-phenyl-2-oxazolidinone as a chiral auxiliary that may be recycled. In general, however, that route is of limited use for scale-up, principally due to the low temperature required for the reactions, the use of pyrophoric reagent, such as, butyl lithium, to side reactions, and due to a low overall yield. [0007] Another process is disclosed in U.S. Patent Application Publication No. 2003/0212290, which discloses asymmetric hydrogenation of a cyano-substituted olefin, compound 7, to produce a cyano precursor of (S)-3-(aminomethyl)-5-methyl hexanoic acid, compound 8, as seen in scheme 2. [0008] Subsequent reduction of the nitrile in compound 8 by catalytic hydrogenation produces (S)-Pregabalin. The cyano hexenoate starting material, compound 7, is prepared from 2-methyl propanal and acrylonitrile (Yamamoto et al, Bull. Chem. Soc. Jap., 58, 3397 (1985)). However, the disclosed method requires carbon monoxide under high pressure, raising serious problems in adapting this scheme for production scale processes. [0009] A process published by G. M. Sammis, et al., J. Am. Chem. Soc., 125(15), 4442-43 (2003), takes advantage of the asymmetric catalysis of cyanide conjugate addition reactions. The method discloses the application of aluminum salen catalysts to the conjugate addition of hydrogen cyanide to α,β-unsaturated imides as shown in scheme 3. Reportedly, TMSCN is a useful source of cyanide that can be used in the place of HCN. Although the reaction is highly selective, this process is not practicable for large scale production due to the use of highly poisonous reagents. Moreover, the last reductive step requires high pressure hydrogen, which only adds to the difficulties required for adapting this scheme for a production scale process. [0010] In 1989, Silverman reported a convenient synthesis of 3-alkyl-4-amino acids compounds in S YNTHESIS , Vol. 12, 953-944 (1989). Using 2-alkenoic esters as a substrate, a series of GABA analogs were produced by Michael addition of nitromethane to α,β-unsaturated compounds, followed by hydrogenation at atmospheric pressure of the nitro compound to amine moiety as depicted in scheme 4. [0011] Further resolution of compound 14 may be employed to resolve Pregabalin. This, of course, results in the loss of 50 percent of the product, a serious disadvantage. However, the disclosed methodology reveals that the nitro compound can serve as an intermediate for the synthesis of 3-alkyl-4-amino acids. [0012] Recent studies have indicated that cinchona alkaloids are broadly effective in chiral organic chemistry. A range of nitroalkenes were reportedly treated with dimethyl or diethyl malonate in THF in the presence of cinchona alkaloids to provide high enantiomeric selectivity of compound 15, and its analogues. For example, see H. Li, et al., J. Am. Chem. Soc., 126(32), 9906-07 (2004). These catalysts are easily accessible from either quinine or quinidine, and are reportedly highly efficient for a synthetically C—C bond forming asymmetric conjugate addition as shown in scheme 5. [0013] R 3 represents several alkyl and aryl groups. The scope of the reaction has been extended to other nitroolefins and applied to prepare ABT-546 employing bis(oxazoline)Mg(OTf) 2 . See, for example, D. M. Barnes, et al., J. Am. Chem. Soc., 124(44), 13097-13105 (2002). [0014] Other groups have investigated a new class of bifunctional catalysts bearing a thiourea moiety and an amino group on a chiral scaffold. See T. Okino, et al., J. Am. Chem. Soc., 127(1), 119-125 (2005). On the basis of a catalytic Michael addition to the nitroolefin with enantiomeric selectivity, they were able to prepare a series of analogues of compound 15. [0015] Thus, there is a need in the art for new processes for the preparation of (S)-Pregabalin that does not suffer from the disadvantages mentioned above. SUMMARY OF THE INVENTION [0016] In one embodiment, the invention encompasses (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-aryl-alkyl]amino}ethyl)hexanoic acid of the following formula 24, wherein Ar is a C 6-10 aromatic group, and R is a straight or branched C 1-4 alkyl, ester or carboxylic acid. [0017] In another embodiment, the invention encompasses a (R)-3-isobutylpentanedioic acid amide-((S)-1-aryl-alkyl)amide of the following formula 25, wherein Ar is a C 6-10 aromatic group, and R is a straight or branched C 1-4 alkyl, ester or carboxylic acid. [0018] In another embodiment, the invention encompasses a (S)-4-methyl-2-{[((S)-1-aryl-alkyl-carbamoyl)-methyl]pentyl}carbamic acid methyl ester of the following formula 26, wherein Ar is a C 6-10 aromatic group, and R is a straight or branched C 1-4 alkyl, ester or carboxylic acid. [0019] In another embodiment, the invention encompasses a (S)-2-carbamoylmethyl-4-methylpentyl)carbamic acid alkyl ester of the following formula 27, wherein R′ is a straight or branched C 1-5 alkyl. [0020] In another embodiment, the invention encompasses a process for preparing (S)-Pregabalin comprising: preparing a compound of the following formula 24 wherein Ar is a C 6-10 aromatic group and R is a straight or branched C 1-4 alkyl, ester, or carboxylic acid; converting the compound of formula 24 into a compound of the following formula 25, wherein Ar is a C 6-10 aromatic group, and R is a straight or branched C 1-4 alkyl, ester, or carboxylic acid; converting the compound of formula 25 into a compound of the following formula 26 wherein Ar is a C 6-10 aromatic group, and R is a straight or branched C 1-4 alkyl, ester, or carboxylic acid; converting the compound of formula 26 into a compound of the following formula 27 wherein R′ is a straight or branched C 1-5 alkyl; and converting the compound of formula 27 into (S)-Pregabalin. [0021] The compound of formula 24 is preferably prepared by a process comprising: combining a chiral amine of the following formula 23, wherein Ar is a C 6-10 aromatic group, and R is a straight or branched C 1-4 alkyl, ester, or carboxylic acid, an organic solvent selected from at least one of C 6-10 aromatic hydrocarbons, substituted aromatic hydrocarbons, C 2-8 ethers, halogenated hydrocarbons, straight or branched C 1-4 alcohols, C 3-8 esters, straight, branched or cyclic C 1-6 alkanes, or C 3-8 ketones, and at least one base, to obtain a mixture; cooling the mixture to a temperature of about −70° C. to about 10° C.; adding 3-isobutyl glutaric anhydride to the mixture; maintaining the mixture at a temperature of about −70° C. to about 10° C. for at least about one hour to obtain the compound of formula 24; and recovering the compound of formula 24 from the mixture. [0022] The compound of formula 24 is preferably converted into the compound of formula 25 by a process comprising: combining at a temperature of about 20° C. to about −30° C., the compound of formula 24 and at least one organic solvent selected from the group consisting of substituted aromatic hydrocarbons, C 6-10 aliphatic hydrocarbons, halogenated carbons, ethers and ketones, an amidation reagent selected from the group consisting of C 1-4 alkyl and C 6-8 aryl haloformates, and acid halides, and a base to form a mixture; maintaining the mixture for about one hour to about two hours at a temperature of about −10° C. to about 20° C.; adding ammonia to obtain the compound of formula 25; and recovering the compound of formula 25 from the mixture. [0023] The compound of formula 25 is preferably converted into the compound of formula 26 by a process comprising: combining a solution of the compound of formula 25 in at least one straight or branched alkyl alcohol, such as methyl, ethyl, isopropyl, n-butyl, isobutyl, or t-butyl alcohol, preferably, methanol or ethanol, at a temperature of about −25° C. to about −45° C. with bromine, in the presence of at least one base, to obtain a basic mixture; warming the basic mixture to a temperature of about 50° C. to about 70° C., preferably, about 55° C. to about 60° C.; warming the basic mixture for about 1 hour to about 4 hours to obtain the compound of formula 26; and recovering the compound of formula 26 from the basic mixture. Preferably, the compound of formula 26 is obtained in a purity of about 90% to about 100% area by HPLC, more preferably, in a purity of about 92% to about 100% area by HPLC, and, most preferably, in a purity of about 95% to about 100% area by HPLC. [0024] Preferably, the base is a metal alkoxide, such as sodium ethoxide, sodium methoxide, potassium methoxide, potassium ethoxide, or potassium tert-butoxide, and is preferably sodium ethoxide or sodium methoxide. Preferably, the compound of formula 26 is recovered by evaporating the solvent from the basic mixture to form a residue and extracting the compound of formula 26 from the residue. Preferably, the compound of formula 26 is extracted with dichloromethane, ethyl acetate, or toluene. The recovered compound of formula 26 is preferably crystallized from an organic solvent selected from at least one of ethers, esters, hydrocarbons, substituted hydrocarbons, and alcohols. Preferably, the organic solvent is diisopropyl ether, ethyl acetate, cyclohexane, dichloromethane, or methanol. [0025] Preferably, the compound of formula 26 is converted into the compound of formula 27 by a process comprising: combining the compound of formula 26 and a mixture of water and an ether to obtain a mixture; combining the mixture with ammonia and an alkali metal at a temperature of about −30° C. to about −60° C., preferably, about −40° C. to about −30° C., to obtain a reaction mixture; maintaining the reaction mixture for about 4 to about 10 hours until the excess of ammonia is evaporated to obtain the compound of formula 27; and, preferably, recovering the compound of formula 27 from the reaction mixture. Preferably, the ether is tetrahydrofuran or dioxane. Preferably, the ammonia is liquid ammonia. Preferably, the alkali metal is lithium or sodium. Preferably, the compound of formula 27 is recovered by extraction, and, more preferably, the compound of formula 27 is crystallized from an ether, such as diisopropyl ether. [0026] Preferably, the compound of formula 27 is converted to (S)-Pregabalin in a process comprising: combining the compound of formula 27 with an acid to obtain a mixture; maintaining the mixture at a temperature of about 60° C. to about 130° C., preferably, about 80° C. to about 110° C., for about 5 to about 30 hours, preferably, for about 18 to about 30 hours, and, more preferably, for about 5 to about 10 hours, to obtain (S)-Pregabalin; and recovering the (S)-Pregabalin from the mixture. Preferably, the acid is a strong mineral acid, such as hydrochloric acid or sulfuric acid. [0027] Preferably, the (S)-Pregabalin is recovered by a process comprising: adjusting the pH of the mixture to about 3 to about 1; extracting a solution of (S)-Pregabalin from the mixture with an alcohol; adjusting the pH of the solution to about 4 to about 7 to induce the precipitation of (S)-Pregabalin; and recovering the precipitated (S)-Pregabalin. Preferably, the (S)-Pregabalin is obtained in a purity of at least about 98% area by HPLC, and, more preferably, about 99% to about 100% area by HPLC. [0028] In another embodiment, the invention encompasses a process for preparing (S)-Pregabalin comprising: combining a compound of the following formula 26, wherein Ar is a C 6-10 aromatic group, and R is a straight or branched C 1-4 alkyl, ester, or carboxylic acid, with an acid to obtain a mixture; maintaining the mixture at a temperature of about 60° C. to about 130° C., for about 3 hours to about 30 hours, to obtain (S)-Pregabalin; and recovering the (S)-Pregabalin from the mixture. [0029] In another embodiment, the invention encompasses a process for recycling 3-isobutyl glutaric anhydride comprising: preparing a compound of formula 24 from 3-isobutyl glutaric anhydride by the process according to claim 29 ; crystallizing the recovered compound of formula 24 from an organic solvent; removing the crystals from the organic solvent; combining the remaining organic solvent with an acid to obtain a first mixture; maintaining the first mixture at a temperature of about 60° C. to about 130° C., to obtain 3-isobutyl glutaric acid; combining the 3-isobutyl glutaric acid with acetic anhydride to obtain a second mixture; heating the second mixture to a temperature of about 125° C. to about 145° C. to obtain 3-isobutyl glutaric anhydride; and recovering the 3-isobutyl glutaric anhydride from the second mixture. [0030] In another embodiment, the invention encompasses (S)-pregabalin having an enantiomeric purity of about 99% to about 100% area by HPLC, preferably of about 99.9% to about 100% area by HPLC. BRIEF DESCRIPTION OF THE FIGURES [0031] FIG. 1 illustrates an IR spectrum of (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylmethyl]amino}ethyl)hexanoic acid of formula 24A. [0032] FIG. 2 illustrates a 1 H-NMR spectrum of (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylmethyl]amino}ethyl)hexanoic acid of formula 24A. [0033] FIG. 3 illustrates a 13 C-NMR spectrum of (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylmethyl]amino}ethyl)hexanoic acid of formula 24A. [0034] FIG. 4 illustrates a powder X-ray diffraction pattern of (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylmethyl]amino}ethyl)hexanoic acid of formula 24A. DETAILED DESCRIPTION OF THE INVENTION [0035] The present invention provides a stereoselective synthesis of (S)-Pregabalin according to the following scheme: This process allows for obtaining (S)-Pregabalin with a relatively high enantiomeric purity. [0036] The invention encompasses (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-aryl-alkyl]amino}ethyl)hexanoic acids of formula 24, wherein Ar is a C 6-10 aromatic group and R is a straight or branched C 1-4 alkyl, ester, or carboxylic acid. [0037] Preferably, the C 6-10 aromatic group is naphthyl, phenyl, substituted phenyl, or substituted naphthyl, more preferably phenyl. Preferably, the substituted phenyl is a phenyl group substituted with at least one of alkoxy, halogen, alkyl, carboxylic acid, or ester. A preferred alkoxy phenyl is methoxyphenyl. Preferred halogenated phenyls are chlorobenzene, bromobenzene, and fluorobenzene. Preferred alkylated phenyls are either toluene or ethylbenzene. Preferably, the carboxylic acid substituent is —COOH, —CH 2 COOH, —CH(CH 3 )COOH or —C(CH 3 ) 2 COOH. Preferably the ester substituent is a methylester, ethylester, isopropylester, n-butylester, isobutyl, or t-butyl derivative of one of the above-listed carboxylic acid substituents. [0038] Preferably, the C 1-4 alkyl is methyl, ethyl, isopropyl, n-butyl, isobutyl or t-butyl, more preferably, methyl. [0039] When Ar is phenyl and R is methyl, the compound of formula 24 is (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylmethyl]amino}ethyl)hexanoic acid 24A which may be characterized by data selected from: a 13 C-NMR (CDCl 3 , 300 MHz) spectrum having carbon chemical shifts at about 21.74, 22.19, 22.66, 24.95, 29.44, 30.89, 36.73, 38.15, 40.55, 43.45, 48.92, 125.41, 126.06, 127.29, 128.57, 143.01, 171.92 and 176.71 ppm; a 1 H-NMR (CDCl 3 , 75 MHz) spectrum having hydrogen chemical shifts at about 0.77, 1.18, 1.38, 1.56, 2.22, 5.03, 6.59-6.62, 7.11-7.22 and 10.88 ppm; and an IR spectrum having peaks at about 3321.39, 2955.91, 1693.33, 1617.43, 1561.07 and 698.24 cm −1 . [0040] The invention also encompasses isolated (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylmethyl]amino}ethyl)hexanoic acid 24A, preferably in crystalline form. The crystalline form of 24A may be characterized by a powder X-ray diffraction pattern having peaks at about 4.3°, 6.9°, 7.2°, and 7.7° 2θ±0.2° 2θ. The crystalline form of 24A may be further characterized by X-ray powder diffraction peaks at about 6.3°, 8.1°, 9.7°, 10.3°, 11.3°, 12.9°, 13.9°, 15.1°, 15.7°, 17.5°, 18.6°, 19.1°, 20.5°, 20.9°, 21.8°, 22.3°, 23.3°, and 23.8° 2θ±0.2° 2θ. Moreover, the crystalline form of 24A may have a melting range of about 95° C. to about 98° C. [0041] The invention also encompasses isolated (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylmethyl]amino}ethyl)hexanoic acid 24A having an optical purity of at least about 80% area by HPLC, preferably of at least about 93% area by HPLC, more preferably, of about 98% to about 100% area by HPLC, most preferably, of about 99% to about 100% area by HPLC. [0042] The invention also encompasses (R)-3-isobutylpentanedioic acid amide-((S)-1-aryl-alkyl)amides of formula 25, wherein Ar is a C 6-10 aromatic group, and R is a straight or branched C 1-4 alkyl, ester, or carboxylic acid. [0043] Where Ar is phenyl and R is methyl, the compound of formula 25 is (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylmethyl)amide 25A. [0044] The invention further encompasses (S)-4-methyl-2-{[((S)-1-aryl-alkyl-carbamoyl)-methyl]pentyl}carbamic acid methyl esters of formula 26, wherein Ar and R are as defined above for formula 24. [0045] When Ar is phenyl and R and is methyl, the compound of formula 26 is (S)-4-methyl-2-[((S)-1-aryl-alkyl-carbamoyl)-methyl]pentyl}carbamic acid methyl ester 26A. [0046] The invention also encompasses (S)-2-carbamoylmethyl-4-methylpentyl)carbamic acid alkyl esters of formula 27, wherein R′ is a straight or branched C 1-5 alkyl, preferably, methyl. [0047] When R′ is methyl, the compound of formula 27 is (S)-2-carbamoylmethyl-4-methylpentyl)carbamic acid methyl ester 27A. [0048] Further, the invention encompasses a process for the preparation of (S)-Pregabalin via the intermediate compound of formula 24. The process comprises preparing the intermediate of formula 24, converting the intermediate compound of formula 24 into the diamide of formula 25, converting the diamide of formula 25 into the chiral carbamate of formula 26, converting the chiral carbamate into the compound of formula 27, and converting the compound of formula 27 into (S)-Pregabalin. Although each of the compounds of formula 24, formula 25, formula 26, and formula 27 can be isolated prior to conversion, isolation of the compounds of formula 26 and formula 27 is not required. Therefore, once the compound of formula 25 has been prepared and isolated, the preparation of the compound of formula 27 from the compound of formula 25 or the compound of formula 26 and the preparation of (S)-Pregabalin from the compound of formula 25 or the compound of formula 26 does not require isolation of any of the intermediate compounds. Thus, once the compound of formula 25 has been prepared and isolated, (S)-Pregabalin can be prepared in a one-pot process without the isolation of either of the compounds of formula 26 or formula 27. [0049] The intermediate compound of formula 24 may be prepared by combining a chiral amine of formula 23, an organic solvent selected from at least one of C 6-10 aromatic hydrocarbons, substituted aromatic hydrocarbons, C 2-8 ethers, halogenated hydrocarbons, straight or branched C 1-4 alcohols, C 3-8 esters, straight, branched or cyclic C 1-6 alkanes, or C 3-8 ketones, and at least one base, to obtain a mixture, cooling the mixture, and adding 3-isobutyl glutaric anhydride of formula 22 to the mixture to obtain the compound of formula 24, which is then recovered. [0050] The 3-isobutyl glutaric anhydride of formula 22 may be prepared according to the process disclosed in U.S. Pat. No. 5,616,793. [0051] The chiral amine of formula 23 is commercially available, and is used as a chiral auxiliary is a primary amine or a chiral amino acid derivative, wherein Ar and R are as defined above for the compound of formula 24. Preferably, the chiral amine of formula 23 is methylbenzylamine, and more preferably the chiral amine of formula 23 is (S)-methylbenzylamine. [0052] Preferably, the aromatic group is toluene. The preferred ether is selected from tert-butyl methyl ether, tetrahydrofuran, diisopropyl ether, and diethyl ether. Preferably, the halogenated hydrocarbon is dichloromethane. Preferred C 1-4 alcohols are isopropyl alcohol, ethanol, methanol, or n-butanol. Preferably, the ester is selected from ethyl acetate, isopropyl acetate, and isobutyl acetate. A preferred straight, branched or cyclic C 1-6 alkane is either hexane or cyclohexane. Preferred ketones are selected from acetone, methyl isobutyl ketone, and methyl ethyl ketone. The more preferred organic solvent is toluene. [0053] Preferably, the base is an organic base selected from diethyl amine, triethyl amine, di-n-propyl amine, di-isopropyl amine, tertbutylamine, morpholine, piperidine, pyridine, and 4-dimethyl aminopyridine. The most preferred base is 4-dimethyl aminopyridine. [0054] Preferably, the mixture is cooled to a temperature of about −70° C. to about 10° C. before adding the 3-isobutyl glutaric anhydride. [0055] Preferably, the mixture is maintained at a temperature of about −70° C. to about 10° C., more preferably of about 0° C. to about −50° C. and most preferably of about −40° C. to −30° C., before recovering the compound of formula 24. Preferably, the mixture is maintained for at least about one hour, more preferably about one hour to about six hours, most preferably, about one hour to about two hours, before recovering the compound of formula 24. [0056] The order of combining the reacting substances when preparing the compound of formula 24 may influence the purity and the yield of the final product. Preferably, the chiral amine of formula 23 is combined with the base, followed by the addition of the 3-isobutylglutaric anhydride of formula 22. [0057] The compound of formula 24 may be recovered by any method known in the art, such as extracting the organic phase with an aqueous basic solution to convert the acidic product to a salt, and acidifying the aqueous phase with a mineral acid to obtain back the acid product. [0058] The compound of formula 24 may optionally be further purified by a crystallization from an organic solvent selected from at least one of esters, nitriles, ethers, C 4-6 straight, branched, or cyclic hydrocarbons, and C 6-10 substituted aromatic hydrocarbons. A preferred ester is ethyl acetate. Preferably, the nitrile is acetonitrile. A preferred ether is methyl t-butyl ether. A preferred C 6-8 substituted aromatic group is either toluene or xylene. Preferred mixtures are that of xylene and ethyl acetate, hexane and ethyl acetate, cyclohexane and ethyl acetate, and toluene and ethyl acetate. The most preferred mixture is that of toluene and ethyl acetate. [0059] The compound of formula 24 obtained by the above-described process has an optical purity of at least about 80% area by HPLC, preferably of at least about 93% area by HPLC, more preferably of about 98% to about 100% area by HPLC, and most preferably of about 99% to about 100% area by HPLC. [0060] The recovered compound of formula 24 is then converted to the diamide of formula 25, in a process comprising combining a mixture of the compound of formula 24 and at least one organic solvent selected from substituted aromatic hydrocarbons, C 6-10 aliphatic hydrocarbons, halogenated carbons, ethers and ketones, an amidation reagent selected from C 1-4 alkyl and C 6-8 aryl haloformates, and acid halides, and at least one base, and adding ammonia to obtain the compound of formula 25, which is then recovered. [0061] Preferably, the ammonia is provided in an aqueous solution, i.e., ammonium hydroxide. [0062] Preferably, the C 1-4 alkyl halo formate is a ethyl or methyl derivative of a chloro or bromoformate. Preferably, the C 6-8 aryl halo formate is a benzyl chloro or bromoformate. Preferred acid halides are acetyl, pivaloyl, oxaloyl or benzoyl chlorides and bromides. The most preferred haloformate is either ethylchloroformate or methylchloroformate. The more preferred acid halide is acetyl, pivaloyl, oxaloyl, or benzoyl chlorides. The most preferred amidation reagent is either ethylchloroformate or methylchloroformate. [0063] Preferably, the substituted aromatic hydrocarbon is either toluene or xylene. A preferred C 6-10 aliphatic hydrocarbon is either hexane or heptane. Preferred ketones are acetone, methyl ethyl ketone, or methyl isobutyl ketone. Preferably, the ether is diethyl ether, diisopropyl ether, or tert-butyl methyl ether. Preferably, the halogenated hydrocarbon is dichloromethane. The more preferred organic solvent is either acetone or dichloromethane. [0064] Preferably, the base is an organic base selected from diethyl amine, triethyl amine, di-n-propyl amine, di-isopropyl amine, tri-n-butyl amine, morpholine, piperidine, pyridine, and 4-dimethyl aminopyridine. The preferred base is either 4-dimethyl aminopyridine or triethyl amine. [0065] Preferably, the mixture of the compound of formula 24 and an organic solvent is combined with the amidation reagent and the base at a temperature of about 20° C. to about −30° C., more preferably, of about −10° C. to about −20° C. Preferably, the compound of formula 24 is compound 24A. [0066] Preferably, the mixture is maintained at a temperature of about −10° C. to about −20° C. before the addition of ammonia. Preferably, the mixture is maintained for about one hour to about two hours before the addition of ammonia. [0067] The compound of formula 25 may be recovered by known methods in the art, such as, filtering and drying. [0068] The compound of formula 25 is obtained by the above process having a purity of at least about 80% area by HPLC, more preferably of at least about 95% area by HPLC. [0069] Then, the recovered compound of formula 25 is reacted with bromine in a Hoffman reaction under basic conditions. The process comprises combining a solution of a compound of formula 25 in at least one straight or branched alkyl alcohol with bromine, in a presence of at least one base, to obtain a basic mixture, and warming the basic mixture to obtain the chiral carbamate of formula 26, which is then recovered. [0070] Preferably the combining step is performed at a temperature of about −25° C. to about −45° C. [0071] Preferably, the base is a metal alkoxide selected from sodium ethoxide, sodium methoxide, potassium methoxide, potassium ethoxide, and potassium tert-butoxide. The more preferred base is either sodium ethoxide or sodium methoxide. [0072] Preferably, the basic mixture is warmed to a temperature of about 50° C. to about 70° C., more preferably to a temperature of about 55° C. to about 60° C. [0073] Preferably, the straight or branched alkyl alcohol is methyl, ethyl, isopropyl, n-butyl, isobutyl, or t-butyl alcohol, more preferably methanol or ethanol. [0074] Preferably, the basic mixture is warmed for about 1 hour to about 4 hours before recovering the compound of formula 26. [0075] The compound of formula 26 may be recovered by evaporating the solvent and further extracting with a solvent selected from dichloromethane, ethylacetate and toluene, followed by drying over a drying agent, such as, magnesium sulfate, followed by evaporating the solvent. [0076] The recovered compound of formula 26 may be purified by crystallization from at least one of an ether, ester, hydrocarbone, substituted hydrocarbon, or alcohol. Preferably, the compound of formula 26 is crystallized from at least one of diisopropyl ether, ethyl acetate, cyclohexane, dichloromethane, or methanol. [0077] The compound of formula 26 is obtained by the above process having a purity of at least about 80% area by HPLC, preferably of about 90% to about 100%, area by HPLC, more preferably of about 92% to about 100% area by HPLC, and most preferably, of about 95% to about 100% area by HPLC. [0078] The amide moiety of the recovered compound of formula 26 is then converted to a primary amide moiety, to give the compound of formula 27 in a process comprising combining the compound of formula 26 and a mixture of water and an ether to obtain a mixture, combining the mixture with ammonia and an alkali metal to obtain a reaction mixture, and evaporating the excess of ammonia to obtain the compound of formula 27. [0079] Preferably, the mixture containing the compound of formula 26 and a mixture of water and ether is combined with ammonia and an alkali metal at a temperature of about −30° C. to about −60° C., more preferably at a temperature of about −40° C. to about −30° C. [0080] Preferably the ether is either tetrahydrofuran or dioxane. [0081] Preferably the ammonia is liquid. [0082] The preferred alkali metal is either lithium or sodium. [0083] Preferably, the excess ammonia is evaporated by maintaining the reaction mixture for about 4 to about 10 hours. [0084] The compound of formula 27 may be recovered by any known method in the art, such as, extraction and drying over anhydrous sodium sulfate. [0085] The compound of formula 27 may optionally be purified by crystallization from ether, preferably, diisopropyl ether. [0086] The compound of formula 27 is then converted to (S)-Pregabalin in a process comprising combining the compound of formula 27 with an acid to obtain a mixture and recovering (S)-Pregablin from the mixture. [0087] Preferably, the acid is a strong mineral acid, more preferably either hydrochloric acid or sulfuric acid. [0088] Preferably, the mixture is maintained at a temperature of about 60° C. to about 130° C., more preferably of about 80° C. to about 110° C., before recovering the (S)-Pregabalin. [0089] Preferably, the mixture is maintained for about 5 to about 30 hours before recovering the (S)-Pregabalin. [0090] Preferably, the mixture is maintained for about 18 to about 30 hours, when the mineral acid is hydrochloric acid and for about 5 to about 10 hours, when the mineral acid is sulfuric acid, before recovering the (S)-Pregabalin. [0091] (S)-Pregabalin may be recovered by adjusting the pH of the mixture to about 3 to about 1, preferably by addition of a strong base; extracting a solution of (S)-Pregabalin from the mixture with an alcohol; adjusting the pH of the solution to about 4 to about 7, preferably with an inorganic or an organic base, to induce the precipitation of (S)-Pregabalin; and recovering the precipitated (S)-Pregabalin. [0092] (S)-Pregabalin obtained by the above process has at least about 80% enantiomeric purity by area HPLC, preferably at least about 93% area by HPLC, more preferably, about 98% to about 100% area by HPLC, even more preferably, about 99% to about 100% area by HPLC, and most preferably of about 99.9% to about 100% area by HPLC. [0093] In an alternative process, the chiral carbamate compound of formula 26 may be converted directly to (S)-Pregabalin. The process comprises combining the compound of formula 26 with an acid to obtain a mixture and maintaining the obtained mixture at a temperature of about 60° C. to about 130° C., for about 3 hours to about 30 hours, to obtain (S)-Pregabalin, which is then recovered. [0094] Preferably, the acid is a strong mineral acid. Preferably, the strong mineral acid is selected from a group consisting of hydrochloric acid, hydrobromic acid, and sulfuric acid. [0095] Preferably, the mixture is maintained at temperature of about 80° C. to about 125° C. [0096] Preferably, the mixture is maintained for about 10 to about 30 hours when the mineral acid is hydrochloric acid, for about 5 to about 10 hours when the mineral acid is sulfuric acid, and for about 3 hours when the mineral acid is hydrobromic acid. [0097] (S)-Pregabalin may be recovered by the same method described for the reaction of converting the compound of formula 27 to (S)-Pregabalin. [0098] Further, 3-isobutyl glutaric anhydride of formula 22 can be regenerated by a process comprising the steps of combining the filtrate obtained from the crystallization process of compound of formula 24A with an acid to form a first mixture, recovering 3-isobutyl glutaric acid of the following formula 28 from the first mixture, combining the 3-isobutyl glutaric acid with acetic anhydride to obtain a second mixture, and recovering 3-isobutyl glutaric anhydride of formula 22 from the second mixture, which may then be reused. [0099] Preferably, the acid is a strong mineral acid, more preferably either a 4N to 12N hydrochloric acid or 20 percent to 80 percent sulfuric acid. [0100] Preferably the first mixture is maintained at a temperature of about 60° C. to about 130° C. before recovering the 3-isobutyl glutaric acid. Preferably, when the mineral acid is hydrochloric acid, the first mixture is maintained at temperature of about 100° C. to about 110° C. before recovering the 3-isobutyl glutaric acid. Preferably, when the mineral acid is sulfuric acid, the first mixture is maintained at a temperature of about 60° C. to about 130° C. before recovering the 3-isobutyl glutaric acid. [0101] Preferably, the second mixture is heated to a temperature of about 125° C. to about 145° C., more preferably, of about 130° C. to about 140° C., before recovering the 3-isobutyl glutaric anhydride. [0102] The 3-isobutyl glutaric anhydride may be recovered by any method known in the art, such as, distilling the excess of acetic anhydride and cooling. [0103] In yet another embodiment, the present invention provides pharmaceutical compositions comprising enantiomerically pure (S)-Pregabalin and at least one pharmaceutically acceptable excipient. Such (S)-Pregabalin has at least about 80% enantiomeric purity, preferably of at least about 93% area by HPLC, more preferably, of about 98% to about 100% area by HPLC, even more preferably, of about 99% to about 100% area by HPLC, and most preferably of about 99.9% to about 100% area by HPLC. Such pharmaceutical composition can be prepared by combining (S)-Pregabalin with one or more excipients or adjuvants. Selection of excipients and the amounts to use may be readily determined by the formulation scientist based upon experience and consideration of standard procedures and reference works in the field. [0104] Diluents increase the bulk of a solid pharmaceutical composition, and may make a pharmaceutical dosage form containing the composition easier for the patient and care giver to handle. Diluents for solid compositions include, for example, microcrystalline cellulose (e.g. Avicel®), microfine cellulose, lactose, starch, pregelitinized starch, calcium carbonate, calcium sulfate, sugar, dextrates, dextrin, dextrose, dibasic calcium phosphate dihydrate, tribasic calcium phosphate, kaolin, magnesium carbonate, magnesium oxide, maltodextrin, mannitol, polymethacrylates (e.g. Eudragit®), potassium chloride, powdered cellulose, sodium chloride, sorbitol, and talc. [0105] Solid pharmaceutical compositions that are compacted into a dosage form, such as a tablet, may include excipients whose functions include helping to bind the active ingredient and other excipients together after compression. Binders for solid pharmaceutical compositions include acacia, alginic acid, carbomer (e.g. carbopol), carboxymethylcellulose sodium, dextrin, ethyl cellulose, gelatin, guar gum, hydrogenated vegetable oil, hydroxyethyl cellulose, hydroxypropyl cellulose (e.g. Klucel®), hydroxypropyl methyl cellulose (e.g. Methocel®), liquid glucose, magnesium aluminum silicate, maltodextrin, methylcellulose, polymethacrylates, povidone (e.g. Kollidon®, Plasdone®), pregelatinized starch, sodium alginate, and starch. [0106] The dissolution rate of a compacted solid pharmaceutical composition in the patient's stomach may be increased by the addition of a disintegrant to the composition. Disintegrants include alginic acid, carboxymethylcellulose calcium, carboxymethylcellulose sodium (e.g. Ac-Di-Sol®, Primellose®), colloidal silicon dioxide, croscarmellose sodium, crospovidone (e.g. Kollidon®, Polyplasdone®), guar gum, magnesium aluminum silicate, methyl cellulose, microcrystalline cellulose, polacrilin potassium, powdered cellulose, pregelatinized starch, sodium alginate, sodium starch glycolate (e.g. Explotab®), and starch. [0107] Glidants can be added to improve the flowability of a non-compacted solid composition and to improve the accuracy of dosing. Excipients that may function as glidants include colloidal silicon dioxide, magnesium trisilicate, powdered cellulose, starch, talc, and tribasic calcium phosphate. [0108] When a dosage form such as a tablet is made by the compaction of a powdered composition, the composition is subjected to pressure from a punch and die. Some excipients and active ingredients have a tendency to adhere to the surfaces of the punch and die, which can cause the product to have pitting and other surface irregularities. A lubricant can be added to the composition to reduce adhesion and ease the release of the product from the die. Lubricants include magnesium stearate, calcium stearate, glyceryl monostearate, glyceryl palmitostearate, hydrogenated castor oil, hydrogenated vegetable oil, mineral oil, polyethylene glycol, sodium benzoate, sodium lauryl sulfate, sodium stearyl fumarate, stearic acid, talc, and zinc stearate. [0109] Flavoring agents and flavor enhancers make the dosage form more palatable to the patient. Common flavoring agents and flavor enhancers for pharmaceutical products that may be included in the composition of the present invention include maltol, vanillin, ethyl vanillin, menthol, citric acid, fumaric acid, ethyl maltol, and tartaric acid. [0110] Solid and liquid compositions may also be died using any pharmaceutically acceptable colorant to improve their appearance and/or facilitate patient identification of the product and unit dosage level. [0111] In liquid pharmaceutical compositions of the present invention, the active ingredient and any other solid excipients are suspended in a liquid carrier such as water, vegetable oil, alcohol, polyethylene glycol, propylene glycol or glycerin. [0112] Liquid pharmaceutical compositions may contain emulsifying agents to disperse uniformly throughout the composition an active ingredient or other excipient that is not soluble in the liquid carrier. Emulsifying agents that may be useful in liquid compositions of the present invention include, for example, gelatin, egg yolk, casein, cholesterol, acacia, tragacanth, chondrus, pectin, methyl cellulose, carbomer, cetostearyl alcohol, and cetyl alcohol. [0113] Liquid pharmaceutical compositions of the present invention may also contain a viscosity enhancing agent to improve the mouth-feel of the product and/or coat the lining of the gastrointestinal tract. Such agents include acacia, alginic acid bentonite, carbomer, carboxymethylcellulose calcium or sodium, cetostearyl alcohol, methyl cellulose, ethylcellulose, gelatin guar gum, hydroxyethyl cellulose, hydroxypropyl cellulose, hydroxypropyl methyl cellulose, maltodextrin, polyvinyl alcohol, povidone, propylene carbonate, propylene glycol alginate, sodium alginate, sodium starch glycolate, starch tragacanth, and xanthan gum. [0114] Sweetening agents such as sorbitol, saccharin, sodium saccharin, sucrose, aspartame, fructose, mannitol, and invert sugar may be added to improve the taste. [0115] Preservatives and chelating agents such as alcohol, sodium benzoate, butylated hydroxy toluene, butylated hydroxyanisole, and ethylenediamine tetraacetic acid may be added at levels safe for ingestion to improve storage stability. [0116] According to the present invention, a liquid composition may also contain a buffer such as gluconic acid, lactic acid, citric acid or acetic acid, sodium gluconate, sodium lactate, sodium citrate, or sodium acetate. [0117] Selection of excipients and the amounts used may be readily determined by the formulation scientist based upon experience and consideration of standard procedures and reference works in the field. [0118] The solid compositions of the present invention include powders, granulates, aggregates, and compacted compositions. The dosages include dosages suitable for oral, buccal, rectal, parenteral (including subcutaneous, intramuscular, and intravenous), inhalant, and ophthalmic administration. Although the most suitable administration in any given case will depend on the nature and severity of the condition being treated, the most preferred route of the present invention is oral. The dosages may be conveniently presented in unit dosage form and prepared by any of the methods well known in the pharmaceutical arts. [0119] Dosage forms include solid dosage forms like tablets, powders, capsules, suppositories, sachets, troches, and losenges, as well as liquid syrups, suspensions, and elixirs. [0120] The dosage form of the present invention may be a capsule containing the composition, preferably a powdered or granulated solid composition of the invention, within either a hard or soft shell. The shell may be made from gelatin, and, optionally, contain a plasticizer such as glycerin and sorbitol, and an opacifying agent or colorant. [0121] The active ingredient and excipients may be formulated into compositions and dosage forms according to methods known in the art. [0122] A composition for tableting or capsule filling can be prepared by wet granulation. In wet granulation, some or all of the active ingredients and excipients in powder form are blended, and then further mixed in the presence of a liquid, typically water, that causes the powders to clump into granules. The granulate is screened and/or milled, dried, and then screened and/or milled to the desired particle size. The granulate may then be tableted or other excipients may be added prior to tableting, such as a glidant and/or a lubricant. [0123] A tableting composition can be prepared conventionally by dry blending. For example, the blended composition of the actives and excipients may be compacted into a slug or a sheet, and then comminuted into compacted granules. The compacted granules may subsequently be compressed into a tablet. [0124] As an alternative to dry granulation, a blended composition may be compressed directly into a compacted dosage form using direct compression techniques. Direct compression produces a more uniform tablet without granules. Excipients that are particularly well suited for direct compression tableting include microcrystalline cellulose, spray dried lactose, dicalcium phosphate dihydrate and colloidal silica. The proper use of these and other excipients in direct compression tableting is known to those in the art with experience and skill in particular formulation challenges of direct compression tableting. [0125] A capsule filling of the present invention may comprise any of the aforementioned blends and granulates that were described with reference to tableting, however, they are not subjected to a final tableting step. [0126] In another embodiment, the present invention provides a method of treating a patient comprising administering to a patient in need thereof a therapeutically effective amount of the above crystalline form of O-desmethylvenlafaxine. Preferably, the patient suffers from a condition which may be treated with a norepinephrine or a serotonin re-uptake inhibitor. Such patient may be suffering from depression. [0127] The following non-limiting examples are merely illustrative of the preferred embodiments of the present invention, and are not to be construed as limiting the invention, the scope of which is defined by the appended claims. EXAMPLES [0128] Chiral HPLC Analysis Instrument: Waters-2487 Column: CHIRAL PACK AD-H, 250 × 4.6 mm, 5 μm Mobile phase: 2% TFA in n-Hexane/Ethanol - 95/5 Flow: 0.5 ml/minute Temperature: 30° C. Wavelength: 210 nm/UV visible spectrophotometer [0129] 1 H-NMR Analysis F2-Acquisition parameters F2-Processing parameters Instrument dpx 300 Probhd 5 mm Dual Z5 SI 32768 Pulprog zg SF 300.13000292 MHz TD 16384 WDW EM Solvent CDCl 3 SSB 0 NS 8 LB 0.50 Hz DS 0 GB 0 SWH 8992.806 Hz PC 1.4 FIDRES 0.548877 Hz AQ 0.9110004 sec RG 32 DW 55.600 μsec DE 4.50 μsec TE 300.0 K D1 5 sec P1 11.35 μsec SFO1 300.1342018 MHz NUC1 1 H PL1 0 dB [0130] 13 C-NMR Analysis F2-Acquisition parameters F2-Processing parameters Instrument dpx 300 Probhd 5 mm Dual Z5 SI 16384 Pulprog zgdc SF 75.4677549 MHz TD 16384 WDW EM Solvent CDCl 3 SSB 0 NS 5043 LB 10.00 Hz DS 0 GB 0 SWH 18832.393 Hz PC 1.4 FIDRES 1.149438 Hz AQ 0.4350452 sec RG 5792.6 DW 26.550 μsec DE 4.50 μsec TE 300.0 K D11 0.03 sec PL12 17.8 Db Cpdprg2 waltz 16 PCPD2 90.00 μsec SFO2 300.1330013 MHz NUC2 1 H PL2 0 dB D1 1 sec P1 9.4 μsec DE 4.5 μsec SFO1 75.4767751 MHz NUC1 13 C PL1 0 dB [0131] IR Analysis KBr pellets Number of sample scans 16 Number of background scans 16 Scanning parameters 4000-500 cm −1 Resolution 4 Sample gain 8 Mirror velocity 0.6329 Aperture 100 [0132] X-Ray Analysis Instrument SIEMENS “ Model: D-5000 Copper radiation 1.5406 A Scanning parameters 2-50° 2θ. Step scan 0.03° Step time 0.5 second Example 1 Preparation of (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylethyl]amino}ethyl)hexanoic acid compound (24) [0133] A three-necked flask equipped with an addition funnel, thermometer pocket, drying tube and a mechanical stirrer, was charged with toluene (400 ml), (S)-(−)-phenylethylamine (142.35 g, 1.1764 mole), and 4-dimethylaminopyridine (0.7176 g, 0.0059 mole). The mixture was cooled to a temperature of −10° C. to −15° C., followed by addition of a solution of 3-isobutyl glutaric anhydride (100 g, 0.59 mole) in toluene (100 ml), over a period of 45-60 minutes, and stirring for additional 1.5-2 hours, at a temperature of −10° C. to −15° C. The mixture was then extracted with 10% aqueous solution of NaOH (500 ml), and the aqueous phase was washed with toluene (1×250 ml). The pH of the aqueous phase was adjusted to 2-2.5 by adding a solution of hydrochloric acid (1-12N). The aqueous phase was further extracted with toluene (1×800 ml) at a temperature of 70-80° C. The toluene layer was washed with 10% sodium chloride solution {700ml) at a temperature of 70-80° C. followed by crystallization to get 125 g (73.0% yield) of a white solid of (3S)-5-methyl-3-(2-oxo-2-{[(1S)-1-phenylethyl]amino}ethyl)hexanoic acid with an optical purity of 99.75%, as measured by chiral HPLC. Example 2 Preparation of (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylethyl]amino}ethyl)hexanoic acid compound (24) [0134] A three-necked flask equipped with an addition funnel, thermometer pocket, drying tube and a mechanical stirrer, was charged with toluene (400 ml), (S)-(−)-phenylethylamine (38.59 g, 0.0.319 mole), and 4-dimethylaminopyridine (0.358 g, 0.0029 mole). The mixture was cooled to a temperature of −40° C. to −50° C., followed by addition of a solution of 3-isobutyl glutaric anhydride (50 g, 0.294 mole) in toluene (100 ml), over a period of 45-60 minutes, and stirring for additional 1.5-2 hours, at a temperature of −40° C. to −50° C. The mixture was then extracted with 3.5-4.0% aqueous solution of NaOH (1000 ml), and the aqueous phase was washed with toluene (1×250 ml). The pH of the aqueous phase was adjusted to 2-2.5 by adding a solution of hydrochloric acid (1-12N). The aqueous phase was further extracted with ethyl acetate (1×300 ml and 1×100 ml), followed by drying the combined ethyl acetates extracts over anhydrous sodium sulphate, and stripping off the solvents to obtain a residue. The residue was crystallized from ethyl acetate and toluene mixture to get 60.7 g (71.0% yield) of a white solid of (3S)-5-methyl-3-(2-oxo-2-{[(1S)-1-phenylethyl]amino}ethyl) hexanoic acid with an optical purity of 99.75%, as measured by chiral HPLC. Example 3 Preparation of (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylethyl]amino}ethyl)hexanoic acid compound (24) [0135] A three-necked flask equipped with an addition funnel, thermometer pocket, drying tube and a mechanical stirrer, was charged with toluene (1000 ml), (S)-(−)-phenylethylamine (266.9 g, 2.206 mole), and 4-dimethylaminopyridine (1.79 g, 0.0147 mole). The mixture was cooled to a temperature of −40° C. to −50° C., followed by addition of a solution of 3-isobutyl glutaric anhydride (250 g, 1.471 mole) in toluene (250 ml), over a period of 45-60 minutes, and stirring for additional 1.5-2 hours, at a temperature of −40° C. to −50° C. The mixture was then extracted with 3.5-4.0% aqueous solution of NaOH (2350 ml), and the aqueous phase was washed with toluene (1×250 ml). The pH of the aqueous phase was adjusted to 2-2.5 by adding a solution of hydrochloric acid (1-12N). The aqueous phase was further extracted with ethyl acetate (1×1250 ml and 1×500 ml), followed by drying the combined ethyl acetates extracts over anhydrous sodium sulphate, and stripping off the solvents to obtain a residue. The residue was crystallized from toluene to get 344 g (80.5% yield) of a white solid of (3S)-5-methyl-3-(2-oxo-2-{[(1S)-1-phenylethyl]amino}ethyl)hexanoic acid with an optical purity of 98.69%, as measured by chiral HPLC. Example 4 Preparation of (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylethyl)amide (25) [0136] A three-necked flask equipped with an addition funnel, thermometer pocket, drying tube and a mechanical stirrer, was charged with methylene dichloride (1000 ml), (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylethyl]amino}ethyl)hexanoic acid compound (24) (200 g, 0.687 mole), and with triethylamine (7.65 g, 0.756 mole), and cooled to 0°-5° C. followed by addition of ethyl chloroformate (90 g, 0.825 mole). The mixture was stirred for 1-2 hours at a temperature of 20° C. to 25° C., followed by quenching with 25% aqueous ammonia (1000 ml). The resulted slurry was filtered and washed with water and dried to get 140 g (70.0% yield) of a white solid of (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylethyl)amide of formula 25A, with a purity of 95%, as measured by HPLC. Example 5 Preparation of (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylethyl)amide (25) [0137] A three-necked flask equipped with an addition funnel, thermometer pocket, drying tube and a mechanical stirrer, was charged with methylene dichloride (500 ml), (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylethyl]amino}ethyl)hexanoic acid compound (24) (100 g, 0.343 mole), and with triethylamine (41.67g, 0.412mole), and cooled to −15° C. to −20° C. followed by addition of ethyl chloroformate (39.1 g, 0.36 mole). The mixture was stirred for 1-2 hours at a temperature of −15° C. to −20° C., followed by quenching over a solution of 20% aqueous ammonia (280 ml). The dichloromethane was distilled out from the mass followed by filtering the resulted slurry, washed with water and dried to get 87 g (87% yield) of a white solid of (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylethyl)amide of formula 25A, with a purity of 98%, as measured by HPLC. Example 6 Preparation of (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylethyl amide (25) [0138] A three-necked flask equipped with an addition funnel, thermometer pocket, drying tube and a mechanical stirrer, was charged with methylene dichloride (125 ml), (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylethyl]amino}ethyl) hexanoic acid compound (24) (25 g, 0.086 mole), triethyl amine (10.43 g, 0.129 mole), and cooled to 0°-5° C. followed by addition of pivaloyl chloride (12.43 g, 0.103 mole). The mixture was stirred for 1-2 hours at a temperature of 20° C. to 25° C., followed by quenching with 20% aqueous ammonia (250 ml). The resulted slurry was filtered and washed with water and dried to get 15.2 g (61% yield) of a white solid of (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylethyl)amide of formula 25A, with a purity of 95%, as measured by HPLC. Example 7 Preparation of (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylethyl) amide (25) [0139] A three-necked flask equipped with an addition funnel, thermometer pocket, drying tube and a mechanical stirrer, was charged with acetone (125 ml), (3S)-5-methyl-3-(2-oxo-2{[(1S)-1-phenylethyl]amino}ethyl) hexanoic acid compound (24) (25 g, 0.086 mole), triethyl amine (10.43 g, 0.129 mole), and cooled to 0-5° C. followed by addition of pivaloyl chloride (12.43 g, 0.103 mole). The mixture was stirred for 1-2 hours at a temperature of 20° C. to 25° C., followed by quenching with 20% aqueous ammonia (250 ml). The resulted slurry was filtered and washed with water and dried to get 10.68 g (43.4% yield) of a white solid of (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylethyl)amide of formula 25, with a purity of 95.4%, as measured by HPLC. Example 8 Preparation of {(S)-4-methyl-2-[((S)-1-phenylethylcarbamoyl)-methyl]pentyl}carbamic acid methyl ester (26) [0140] A three-necked flask equipped with an addition funnel, thermometer pocket, drying tube and a mechanical stirrer, was charged with methanol (1400 ml), and cooled to −40° to −45° C. followed by addition of sodium methoxide (130 g, 2.413 mole). A solution of bromine (154.48 g, 0.965 mole) in methanol (300 ml) was slowly added at about −40 to −45° C. followed by addition of (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylethyl)amide, of formula 25 (140 g, 0.48 mole), in methanol (560 ml). The mixture was gradually warmed to a temperature of 0° C. and then to 55-60° C., followed by stirring for 2 to 3 hours. The solvent was then stripped off and water was added to the mass. The resulted slurry was further extracted with methylene dichloride (1×500 ml and 1×250 ml), followed by drying the combined methylene dichloride extracts over anhydrous sodium sulphate, and stripping off the solvents to obtain a residue. The residue was crystallized from diisopropyl ether to get 115 g (74.2.0% yield) of a white solid of {(S)-4-methyl-2-[((S)-1-phenylethylcarbamoyl)-methyl]pentyl}carbamic acid methyl ester (26) with a purity of 92%, as measured by HPLC. Example 9 Preparation of {(S)-4-methyl-2-[((S)-1-phenylethylcarbamoyl)-methyl]pentyl}carbamic acid methyl ester (26) [0141] A three-necked flask equipped with an addition funnel, thermometer pocket, drying tube and a mechanical stirrer, was charged with methanol (2000 ml), and cooled to −15° to −20° C., (R)-3-isobutylpentanedioic acid amide-((S)-1-phenylethyl)amide, of formula 25 (100 g, 0.344 mole) followed by addition of sodium methoxide (74.5 g, 1.38 mole). Bromine (82.56 g, 0.516 mole) was slowly added at about −15 to −20° C. The mixture was gradually warmed to a temperature of 0° C. and then to 55-60° C., followed by stirring for 2 to 3 hours. The solvent was then stripped off and water was added to the mass. The resulted slurry was further extracted with methylene dichloride (1×500 ml), followed by washing the methylene dichloride extract with water and brine solution. The solvent was stripped off and the residue was crystallized from a mixture of methylene dichloride and cyclohexane to get 95 g (85.8.0% yield) of a white solid of {(S)-4-methyl-2-[((S)-1-phenylethylcarbamoyl)-methyl]pentyl}carbamic acid methyl ester (26) with a purity of 93%, as measured by HPLC. Example 10 Preparation of {(S)-2-carbamoylmethyl-4-methylpentyl)carbamic acid methyl ester (27) [0142] A 2 l, four-necked flask, equipped with a mechanical stirrer, thermometer pocket and a liquid ammonia inlet, was charged with {(S)-4-methyl-2-[((S)-1-phenylethylcarbamoyl)-methyl]pentyl}carbamic acid methyl ester (26) (25 g, 0.078 mole), tetrahydrofuran (175 ml), and water (25 ml). The reaction mixture was cooled to −40° to −60° C. and liquid ammonia (1000 ml) was added followed by addition of small pieces of sodium metal (7.2 g). The resultant reaction mixture was stirred vigorously for 4-10 hours until the ammonia had evaporated. Water (100 ml) was added to the reaction mass under N 2 atmosphere at 5°-10° C., followed by separating the phases. The organic layer was separated and dried over anhydrous sodium sulphate and the solvent was stripped off. The residue was crystallized from diisopropyl ether to get 10.2 g (60% yield) of {(S)-2-carbamoylmethyl-4-methylpentyl)carbamic acid methyl ester with purity of 73% as measured by HPLC. Example 11 Regeneration of 3-isobutylglutaric acid [0143] A 3 l four-necked flask, equipped with a mechanical stirrer, thermometer pocket, and condenser, was charged with a residue after crystallization of compound 24A, (250 g) from examples 1 and 2, and 70% sulfuric acid (2500 g). The reaction mixture was refluxed at 115°-125° C. for 5-10 hours, and then cooled to 20°-25° C. and diluted with water. The aqueous layer was extracted with toluene (1×1000 ml and 1×500 ml). The combined organic phase was extracted with 5% sodium hydroxide solution (1500 ml), and the pH of the aqueous layer was adjusted to 1.5-2 with concentrated hydrochloric acid, followed by extractions with toluene (1×600 ml and 1×400 ml). The combined organic layers were dried over anhydrous sodium sulphate and the solvent was stripped off to obtain 3-isobutyl glutaric acid (128 g) in purity of 94% as measured by GC. [0144] 3-isobutylglutaric acid is characterized by: 1. IR (KBr): 1713.27 cm −1 . 2. 1 HNMR (CDCl3): δ 0.89-0.92 (d, 6H), 1.25 (t, 2H), 1.6-1.69 (septet, 1H), 2.42 (s, 4H), 11.96 (s,2H). 3. 13 C NMR (CDCl 3 ): δ 22.39, 25.06, 28.11, 29.50, 38.45, 43.38, 179.17, 203. Example 12 Converting 3-isobutylglutaric acid to 3-isobutylglutaric anhydride, Compound 22 [0148] A 1 l, four-necked flask equipped with a mechanical stirrer, thermometer pocket and condenser, was charged with 3-isobutyl glutaric acid (500 g) and acetic anhydride (326 g). The reaction mixture was refluxed at 135°-1450° C. for 2.5-3 hours, followed by distilling out the unreacted acetic anhydride at 147°-155° C., and then the distillation was continued under vacuum to ensure removal of traces of unreacted acetic anhydride. The residue was cooled to 25°-30° C. to obtain 445 g of 3-isobutylglutaric anhydride. Example 13 Preparation of (S)-Pregabalin [0149] A 0.2 l reactor was loaded with 6N hydrochloric acid (100 ml) containing compound 27 (12 g, 0.055 mole), and was heated to 100°-110° C. for 12-24 hours, and then cooled to room temperature, i.e., about 20° to about 25° C. An aqueous 40% sodium hydroxide solution was added in an amount sufficient to provide a pH of 1. The solution was then extracted with 37 ml of iso-butanol, the organic layer was separated, and Bu 3 N was added in an amount sufficient to provide a pH of 4. The (S)-Pregabalin was precipitated, filtered, and washed with 10 ml of iso-butanol. After drying at 55° C. under vacuum, (S)-Pregabalin was obtained as white crystals in a 21.5% yield. Purity: 99.9% area by HPLC. Example 14 Preparation of (S)-Pregabalin [0150] A 0.2 l reactor was loaded with 70% sulfuric acid (200 g) containing compound 26 (10 g, 0.031 mole), and was heated to 115-120° C. for 5-10 hours, and then cooled to room temperature, i.e., about 20° to about 25° C. An aqueous 40% sodium hydroxide solution was added in an amount sufficient to provide a pH of 1. The solution was then extracted with 35 ml of iso-butanol, the organic layer was separated, and Bu 3 N was added in an amount sufficient to provide a pH of 4. The (S)-Pregabalin was precipitated, filtered, and washed with 10 ml of iso-butanol. After drying at 55° C. under vacuum, (S)-Pregabalin was obtained as white crystals in a 40.4% yield. Purity: 99.95% area by HPLC. Example 15 Preparation of (S)-Pregabalin [0151] A 0.2 l reactor was loaded with 70% sulfuric acid (200 g) containing compound 26 (10 g, 0.031 mole), and was heated to 115-120° C. for 5-10 hours, and then cooled to room temperature, i.e., about 20° to about 25° C. An aqueous 40% sodium hydroxide solution was added in an amount sufficient to provide a pH of 1. The solution was then extracted with 50 ml of isopropanol, the organic layer was separated, and NH 4 OH was added in an amount sufficient to provide a pH of 4. The (S)-Pregabalin was precipitated, filtered, and washed with 10 ml of isobutanol. After drying at 55° C. under vacuum, (S)-Pregabalin was obtained as white crystals in a 50.4% yield. Purity: 99.05% area by HPLC. Example 16 Preparation of (S)-Pregabalin [0152] A flask was loaded with 47% HBr (12 ml), water (6 ml), and compound 26 (6 g), and then was heated to reflux for 3 hours. The solution was cooled to room temperature, and water (12 ml) was added. An aqueous 47% sodium hydroxide solution was added to obtain pH of 3. The solution was then extracted twice with isobutanol (15 ml), the combined organic layers were evaporated and fresh isobutanol was added (15 ml). Bu 3 N (3.8 g) was added. The mixture was cooled to 2° C. for 1 hour, then (S)-Pregabalin was filtered, and washed with of iso-butanol (3 ml). After drying at 55° C. under vacuum, (S)-Pregabalin was obtained as white crystals in a 90% yield. Example 17 Preparation of (S)-Pregabalin [0153] A flask was loaded with 47% HBr (30 ml), water (15 ml), and compound 26 (15 g), and then was heated to reflux for 3 hours. The solution was cooled to room temperature and water (30 ml) was added. An aqueous 47% sodium hydroxide solution was added to obtain pH of 3. The solution was then extracted twice with iso-butanol (37.5 ml). The organic layers were combined and Bu 3 N (9.5 g) was added. The mixture was cooled to 2° C. for 1 hour, then (S)-Pregabalin was filtered, and washed with of iso-butanol (10 ml). After drying at 55° C. under vacuum, (S)-Pregabalin was obtained as white crystals in a 51% yield. [0154] While it is apparent that the invention disclosed herein is well calculated to fulfill the objects stated above, it will be appreciated that numerous modifications and embodiments may be devised by those skilled in the art. Therefore, it is intended that the appended claims cover all such modifications and embodiments as falling within the true spirit and scope of the present invention.

The invention encompasses processes for the synthesis of (S)-(+)-3-(aminomethyl)-5-methylhexanoic acid, (S)-Pregabalin, and intermediates of (S)-Pregabalin.

This is a continuation of co-pending application Ser. No. 015,203 filed on 02/17/87, now abandoned. BACKGROUND OF THE INVENTION 1. Field of the Invention The invention relates to the field of magnetoresistive sensors and more particularly to magnetoresistive heads for magnetic disk drives. 2. Brief Summary of the Prior Art Magnetoresistive sensors responsive to a change in resistivity caused by the presence of magnetic fields are increasingly employed as read transducers in the heads of magnetic disk drives primarily because the change of resistivity is independent of disk speed, depending only on the magnetic flux and secondarily because sensor output may be scaled by the sense current. These sensors typically comprise a thin strip of NiFe alloy (Permalloy) magnetized along an easy axis of low coercivity. Many other ferromagnetic alloys are also candidates. The strips are usually mounted in the head such that the easy axis is transverse the direction of disk rotation and parallel to the plane of the disk. The magnetic flux from the disk causes rotation of the magnetization vector of the strip, which in turn causes a change in resistivity to a sense current flowing between lateral contacts. The resistivity varies approximately according to the cosine-squared of the angle between the magnetization vector and the current vector (i.e., delta-rho=rho-max * cosine-squared theta, where theta is the angle between the magnetization and current vectors and rho is the resistivity). Due to this cosine-squared relationship, if the magnetization and current vectors are initially aligned, the initial change in resistivity due to disk magnetic flux is low and undirectional. Typically, therefore, either the easy axis magnetization vector or the current vector is biased to approximately 45° to increase responsiveness to angular change in the magnetization vector and to linearize the sensor output. One problem encountered with magnetoresistive sensors is Barkhausen noise caused by the irreversible motion of magnetic domains in the presence of an applied field, i.e., coherent rotation of the magnetization vector is non-uniform and suppressed, and depends upon domain wall behavior. This noise mechanism is eliminated by creating a single magnetic domain in the sense current region of the strip. Many different means have been employed to both linearize the sensor output and to provide for a single domain in the sense region. An increase in the length of a strip of magnetoresistive material relative to its height is known to contribute to formation of a single magnetic domain in the central portions of the strip. These migrate toward the center under the influence of external fields. However, long strips may be subject to cross-talk in lateral portions of the strip and may conduct magnetic flux from adjacent tracks to the sense region of the strip. Short strips, in contrast, almost invariably spontaneously "fracture" into multiple domains. Efforts have been made to provide single domains in the sensor region by shaping the strip so as to reduce edge demagnetizing fields while providing a relatively short physical dimension in the sensor region. See e.g., Kawakami et al. U.S. Pat. No. 4,503,394, at FIG. 4a, wherein upper and lower horizontal sections with opposed easy axes are connected at the ends with vertical sections to comprise an endless loop. See also, U.S. Pat. No. 4,555,740 in which the strip has two intermediate, upwardly extending legs. However, even shaped strips "fracture" into multiple domains in the presence of strong transverse magnetic fields caused by the inductive write poles between which the magnetoresistive sensors are conventionally mounted (the poles act as soft-magnetic shields to isolate the sensor from magnetic fields not directly adjacent to the sensor). Efforts have also been made to form single domains by providing a longitudinal magnetic field in "long" or shaped strips, prior to reading. Such a magnetic field has to be strong enough to cause the formation of a relatively stable, single domain in the central sensor region. This initialization field is generally provided by a barber pole, which is also used to cant the direction of the sense current relative to the easy axis magnetic vector. For short strips, efforts have been made to maintain single domains by permanent longitudinal biasing from adjacent permanent magnets or atomically coupled antiferromagnetic material which results in exchange biasing. Such biasing means are also provided in some applications to transverse-bias the magnetic vector away from the easy axis to linearize the sensor output, as mentioned above. Both of these biasing schemes (initialization and permanent) have the drawback in that the biasing magnetic field could adversely affect the information prerecorded on the magnetic disk, and further, a permanent biasing field (both transverse and longitudinal) increases the effective anisotropy of the sensor thereby decreasing sensitivity to disk magnetic flux. The barber pole (canted current) design has the additional disadvantage that the effective length of the sensor area is less than the longitudinal distance between the sensor contacts. The barber pole design also requires precise lithographic processes to apply the canted contacts and shorting stripes. Exchange-biasing is not commonly used in practice because of the presence of two dissimilar materials (the magnetoresistive material and the antiferromagnetic material) at an exposed interface. This can lead to corrosion which can destroy a head. Further, because exchange biasing is a quantum-mechanical ineraction effect, reliable atomic interaction is a must, but such processing is difficult and yields are low. Further, the effect has a strong temperature dependence, being substantially reduced in the typical operating environments of conventional disk drives. SUMMARY OF THE INVENTION The invention comprises a series of incremental improvements which address the several problems of prior magnetoresistive heads and leads either alone, or in combination, to an improved head. This disclosure is common to several co-pending applications, each claiming individual improvements which combine to form an improved magnetoresistive sensor and head. These improvements include shaping the strip in the form of a pseudo-ellipse. This shape has a very stable single domain in the central sense region of the strip. Next exchange biasing antiferromagnetic material may be atomically coupled to the ends of an arbitrarily-shaped strip for the purpose of maintaining the central region in a single domain state. Due to the quantum-mechanical effect of the exchange material, the material does not have to cover the entire ends of the strip, but may be recessed away from the exposed interface region reducing the susceptability to corrosion. Once stability has been established via the pseudo ellipsoid shape and/or boundary control exchange stabilization, only two canted contacts are needed to change current direction for the purpose of linearzing the MR sensor. This completely eliminates the need for any barber poles for stabilizing a domain state and reduces the number of electrical contacts to only two--the sense contacts--because barber poles are no longer needed. The canted current design is further improved by patterning the strip to cant the easy axis of the strip relative to the horizontal plane of the magnetic disk and correspondingly relaxing the angular cant of the contacts. This leads a greater effective longitudinal sense region. Further, transverse biasing may be eliminated entirely in coded digital applications where the location of data rather than its magnetic strength or direction is important by operating the sensor in its non-linear mode. While reducing the dynamic range to a small extent, zero crossing determinations from the derivative of the sensed read signals are improved by the increased slope of the non-linear response. Finally, the sensor is preferably located to outside of the inductive write gap to avoid the deleterious effect of multiple domain formations caused by the strong magnetic fields present during write operations. An additional gap structure is added having a broad central shield/pole to shield an elongated magnetoresistive sensor while providing good write/read characteristics. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is an elevation of a pseudo-ellipsoid magnetoresistive sensor strip. FIG. 2 is the strip of FIG. 1 having exchange-biasing material at the ends. FIG. 3 illustrated the essential steps for depositing exchange-biased only on the ends of a magnetoresistive strip. FIG. 4 is an elongated magnetoresistive strip having upwardly projecting ends with exchange-biased material at the ends. FIG. 5 is a cross section of a magnetoresistive head having a magnetoresistive sensor with recessed exchange-biased material. FIG. 6 shows the layer structure of a double-gap magnetoresistive head. FIG. 7 is an elevation view of the essential elements of a double-gap magnetoresistive head. FIG. 8 shows the prior art canted current contacts and electrical circuits connected thereto. FIG. 9 shows the easy axis pattern-biased strip of the present invention and relaxed canted current contacts. FIG. 10 shows an easy axis, patterned-biased pseudo-ellipsoid magnetoresistive strip. FIG. 11 shows the relative response of a magnetoresistive sensor in linear and non-linear modes. FIG. 12 is a pseudo-ellipsoid magnetoresistive sensor having uncanted contacts for a non-linear response. DESCRIPTION OF THE PREFERRED EMBODIMENT FIG. 1 shows the pseudo ellipsoid structure of a magnetoresistive sensor 10 magnetized along an easy axis M. The central portion, indicated by L, has relatively flat sides, rather than curved as in a true ellipse. The aspect ratio, AR of overall length to height is less than 3, but can be greater with no loss of effect. From the central region L, the sides converge to apices in which small magnetic domains 12 and 14 spontaneously form. Preferably W≦L and E, the length of an end, is on the order of L at a minimum, having no known maximum. The structure forms a very stable central region single domain indicated by the large right arrow. Experimentation with this structure shows that a thin layer of 200-500 angstroms NI: 82 Fe: 18 alloy, with an overall length of 25 microns, an L portion of 9 microns, and a width W of 8 microns requires 35 Oe to switch the magnetization vector of the central region to the hard axis, while only 0.75 Oe is required in an unpatterned bulk film. This translates to a factor of 46 improvement. Where high transverse fields can be expected, such as when an unshielded sensor is placed between or next to the poles an inductive write head, longitudinal biasing is still required to initialize or maintain a single domain state. As hitherto discussed, there are many different means for accomplishing this. For example, barber pole biasing generates a longitudinal field. In addition permanent magnetic biasing or exchange biasing can also provide a longitudinal field. A novel stabilization means is disclosed in FIG. 2. Prior exchange stabilization/biasing techniques have been typically prepared by first depositing a ferromagnetic layer upon a substrate and then depositing an antiferromagnetic layer upon the ferromagnetic layer such that after patterning, the two layers coincide. Exchange biasing can result in a signal loss due to shunting effects. Longitudinal exchange fields have a negative temperature dependence. And finally, the possibility exists of corrosion due to bimetallic film structure. The domain stabilization process can be understood by recognizing that if the magnetization is somehow pinned at the boundaries of a thin film strip then the equilibrium magnetization direction can be controlled in the central region between the boundaries. By depositing FeMn in the hatched regions shown in FIG. 2, the previously indicated drawbacks to standard exchange-biasing techniques can be avoided. First, since there is no exchange material in the central active region, there can be no signal loss due to current shunting. Secondly, this stabilization technique is extremely temperature insensitive since the requirement is made that only the direction of the magnetization be fixed; not the magnitude of the longitudinal exchange field. And finally, by proper patterning, the bimetallic interface can be eliminated at any exposed edge. In the preferred embodiment, the exchange biasing material is FeMn because it is electrically conductive. An embodiment having a stable single-domain central region employing exchanged-biased ends is shown in FIG. 2. Here, the strip has a C shape with a relatively narrow central region and with the lateral ends having upwardly extending legs 26 and 28 for conducting the demagnetizing field further away from the central region. This improves the stability of single domains in the central region. Exchange bias material 32, 34 and contact metallization (not shown in the figure) are applied to these ends using the process next described in the pattern shown in the Figure, which generally conforms to the pattern of canted current end contacts (not shown) to be applied later. This pattern of exchange material eliminates edge and end domains and provides a stable central single domain sense region. To avoid the aforementioned corrosion problem at exposed interfaces, the resist pattern is so shaped as to provide a recess S of between the exchange material and the lower edge of the strip 10, the edge exposed to the magnetic disk in most designs. The process for forming the structure shown in FIG. 2 is shown in FIG. 3. Step 1, a strip of magnetoresistive material is evaporated, sputtered or the like on a substrate, not shown for the purpose of clarity, in a uniform magnetic field directed along an easy axis and patterned. Step 2: A photo resist layer is laid down and patterned using conventional processes to form an island resist layer 20 with inwardly inclined sides. Step 3: Next the exchange material 22 is evaporated, sputtered or the like onto the combination. Step 4: deposit contact metallization 23. Step 5: Using a lift-off process, the resist, exchange material 24 and metal 23 clinging to it are removed. FIG. 4 shows a pseudo-ellipsoid strip with exchange material on its ends extending to the flattened central region L. A similar recess S 36 should be provided. FIG. 5 shows a cross section of the exchange-biased strip 10 of either FIG. 2 or 4 mounted between shields 42 and 44 of a typical head above a magnetic disk 50. In the Figure, the exchange material 32 is recessed a short distance 36 above the head surface and the contact metal 38 has a leg 40 extending to the strip 10 and shielding the exchange material 32 from exposure. At least one of the shields 42 and 44 also comprise one of the poles of an inductive write apparatus. The shields are separated via spacers 52 typically composed of a non magnetic material such as Al 2 O 3 . By providing the recess 36, the contact 38 has a leg 40 directly contacting the magnetoresistive material 10. This shields the exchange material 32 from exposure. Most heads land on the disk surface 50 when disk rotation stops, burnishing small amounts of head material away. The amount of recess verses the degree of burnishing determines the lifetime of the head until the exchange material is exposed leading to potential corrosion. The existence of strong transverse magnetic fields causes relatively stable single domain regions to "fracture" into multiple domains, the source of Barkhausen noise. Strong magnetic fields are present between the pole tips of an inductive write apparatus, the conventional location of most magnetoresistive heads. To reduce affect of the inductive write pole tips on the magnetoresistive head, it is known to place the head along side the inductive write pole tips. See e.g., Lee, U.S. Pat. No. 4,321,641. This type of structure requires a soft-magnetic shield, a shield/trailing pole tip, and a leading pole tip. The design of this patent is not entirely satisfactory, primarily due to the extension of the MR material 76, 78 (see FIG. 4 or 7 of '641) beyond the shield of the trailing pole tip 90. The design of FIGS. 6 and 7 provide a very magnetically quiet region for the magnetoresistive sensor. The residual flux from the magnetic poles of the inductive write transducer is so low as to permit the reliable operation of very stable, shaped, single-domain sensors (e.g., the pseudo-ellipse 10 of FIG. 1) without longitudinal biasing. FIG. 6 is a cross section of the essential elements of the improved design. A layer of oxide, preferably aluminum oxide 62 is deposited upon a soft-magnetic substrate 60, preferably NiZn. Next, the magnetoresistive sensor material 64 is deposited in a magnetic field and patterned. (Exchange biasing material may then be deposited and patterned if desired). Metal contacts 66 are then deposited on the magnetoresistive strip 64. A second layer of oxide 68 is then deposited. These two oxide layers, 62 and 68, comprise the read gap. A layer of polyimide or photo resist 70 is then deposited and patterned as shown to remove the layer adjacent the gap end of the head. Next a layer of ferromagnetic material 70 is laid down, preferably NiFe (Permalloy). This layer 70 comprises the trailing pole/shield. Next, a write gap oxide, 75 (aluminum oxide or silicon dioxide), is deposited followed by a second layer of polyimide or photo resist 74. Metal coils 78 are deposited and patterned. Two layers of polyimide or photo-resist 76 are deposited and patterned to remove the portion not adjacent the coils 78. And finally, a final layer of a ferromagnetic material 79 is deposited to encompass the coils and contact the other ferromagnetic layer 72 to form a continuous flux path. After the package is formed, it is typically sealed in a suitable non magnetic material and the gap end processed (usually lapped) to expose the gap(s) and provide a reliable gap height. FIG. 7 is an end view of the essential elements of the double gap head of the preferred embodiment. Spacing layers are omitted for clarity. Shown in the Figure are the ferrite substrate 60, the magnetoresistive strip 64, its lateral metal contacts 66 defining a central sensor region 65 of length L, and the ferromagnetic trailing pole/shield 72, and the leading pole 79. The length of the leading pole 79 defines the written track width via magnetic mirroring with the trailing pole/shield as shown in the Figure. This length corresponds to the length L (plus a processing guard band wherein the length L is intentionally made smaller than the written track width to avoid possible cross talk) of the central region 65 of the magnetoresistive strip 64. Typically, the magnetoresistive strip is longer than a track width to assist in providing a stable central region single domain. It is essential that the trailing pole/shield 72 be as long as the magnetoresistive sensor 64 to completely shield it from side fringing fields originating during the writing process. This makes the lead and trailing poles 79, 72 of different lengths. But it has been discovered that this does not affect the written track width, which is defined by the length of the leading pole 79 and the above mentioned mirroring effects. For many applications, such as audio, linear operation of a magnetoresistive sensor is desirable. As mentioned above, linearization either requires the canting of the easy axis magnetization vector or the canting of the current vector. Canting of the magnetization vector typically increases anisotropy and reduces the range of resistivity change and thus sensitivity of the sensor. Canting the current likewise causes a comparable loss in sensitivity as better illustrated in FIG. 8. FIG. 8 shows a typical canted current biasing technique wherein conductors 80, 82, intimately contacting a magnetoresistive strip 92 of length L provide a canted current from a source 88 generally in the direction L eff between the contacts. The current direction is generally perpendicular to the surfaces 84, 86 of the contacts. These surfaces are generally canted at an angle theta sub B of between 40 and 45 degrees for greatest linearity and sensitivity. The change in resistivity is sensed by means 90 (which can be a voltage sensor if 88 is a constant current source, a transimpedance current sensor if 88 is a constant voltage source, or a power sensor if 88 is a "soft" source). The change is resistivity is generally proportional to the length L eff , which by inspection, is less than the length L between the contacts in the longitudinal direction. L, in turn, is approximately the track width of narrow track and defines the length of the sense region. Thus, the sensitivity of the device is reduced by the ratio L eff /L. Making L eff comparable to track width is not desirable because L would then be long enough to pick up significant cross talk from adjacent tracks. FIG. 9 shows the improved canted current sensor which relaxes the cant of the contact surfaces 84, 86 to an angle theta sub B' of approximately 50 degrees. This substantially increases L eff and thereby sensitivity while maintaining the approximately 40 to 45 degree angle with the easy axis. The reason for this is that the magnetoresistive strip is patterned such that its easy axis is itself canted by an angle theta sub EA of approximately 10 degrees. In the Figure, the contact surfaces 84, 86 are each canted at an angle theta sub B' which is preferably 50 degrees. The magnetoresistive strip's lower edge 96 is parallel to the magnetic disk surface as in the prior art, but the upper edge 98 is patterned at an angle theta sub P thereto to give a resultant easy axis magnetization vector M at an angle theta sub EA of approximately 10 degrees with the lower edge. The strip 94 is formed from a bulk film deposited on a suitable substrate in an uniform magnetic field directed parallel to the lower edge 96. The bulk film is thereafter patterned using conventional lithographic techniques to form the pattern illustrated with the upper edge forming an upwardly extending angle with the lower edge. This shape inherently cants the easy axis magnetization vector upwards, though to a lesser degree than the angle of the upper edge. To achieve the net easy axis rotation of 10 degrees, the designer will have to balance the strength of the undeflected easy axis vector with the size, length, thickness, and composition of the magnetoresistive material with the degree of upward edge angle. In the preferred embodiment, the strip 94 is composed of 80:20 NiFe alloy, approximately 500 angstroms thick, with L approximately equal to 9 microns, h (the height of the sensor at point 104) approximately equal to 8 microns, and theta sub P is 10 degrees, where theta sub B' is 50 degrees and theta sub EA is 10 degrees. FIG. 10 shows an example of a canted easy axis psuedo-ellipsoid 100 and its relative orientation with contact surfaces 84 and 86 (the balance of the contacts are not shown in the Figure. for most digital applications, data is written on the disk in code (e.g., variable length 2,7) where only the location of a transition (pulse peak) rather than its direction and magnitude is important. Pulse amplitude serves the function of triggering qualifiers to discern between signal and noise. Thus except for improved initial sensitivity to magnetization vector rotation, there is no good reason to operate the sensor linearly. Thus the final improvements to the magnetoresistive sensor are to provide no transverse biasing at all, except perhaps the patterned biasing just described, operate the sensor in a non linear mode and design the magnetoresistive sensor and disk flux such that the magnetization vector rotation in response to disk flux is on the order of 40 to 50 degrees. Because the location of a transition (pulse peak) is important, the signal from the disk is conventionally differentiated and the zero crossing detected. Noise makes the location of the zero crossing uncertain, and for this reason, noise ultimately limits the data density. However, by not biasing the sensor, the sensor will operate in its non linear mode (see the equation in the Description of the Prior Art) and the differential will have a steeper zero crossing slope than that of a linearly biased sensor. This increased zero-crossing slope leads to less sensitivity to noise, and to a more accurate detection of zero crossing locations, all other things being equal. In order to achieve a proper non linear signal out of the sensor, the magnetization vector must be made to rotate to a greater degree than in the biased case, see FIG. 11 for an illustration of the principle. The upper portion of the Figure graphs half of the normalized magnetoresistive response (the cosine-squared equation previously mentioned). The lower portions of the Figure graph two disk flux input signals, the left 104 represent the input signal to a non linear magnetoresistive sensor, the right the input to a linear magnetoresistive sensor. While the two signals are shown with greatly different magnitudes, they may actually be the same magnitude if the relative response of the magnetoresistive sensor is proportional to the relative difference shown. Actually, it may be preferable to adjust the relative response of both the disk and the sensor. In a linear mode of operation, the input pulse 106 passes through states 1, 2, 3, and 4 and the sensor responds by moving through resistivity states 1', 2', 3' and 4' (for an oppositely polarized pulse, the states would be on the opposite side of 1'). For all states, the inputs and outputs are linear replications. In a non linear mode, the input signal 104 passes through states A→F and the sensor responds with states A'→F' (an oppositely polarized signal pulse would result in the same output, but from the other half of the resistivity curve.) The output is non linear until region D'→F', where it again becomes a linear replication of the input. It can be seen from the Figure that the total response of the non linear sensor (from A' to F') is greater than the total response from the linear sensor (from 1' to 4'). Thus total sensitivity is greater, and transition centers (pulse peaks) can be more accurately located. Actual sensor output is increase by 25 to 30%. While there are many possible choices of materials to achieve the response indicated in FIG. 11, the preferred choice is a sensor comprised of permalloy and a magnetic disk material with sufficient flux to yield the magnetization vector rotation indicated with a head mounted on a conventional flyer. FIG. 12 shows the preferred magnetoresistive sensor comprised of a pseudo-ellipsoid 10, uncanted contacts 110 and 112, a constant current source 88, a voltage sensor 90. This sensor is preferably mounted in the double gap head of FIGS. 6 and 7. By providing no biasing whatever, the sensor operates in a non linear mode. Its shape and location in the shielded second gap of the double gap head maintains the sensor in a single domain state. If this implementation is not robust enough for a given application, then stability can be further enhanced by providing for exchange material in contact regions 110 and 112 as previously described.

An elongated magnetoresistive sensor strip longitudinally biased to maintain a single domain sense region by exchange-biasing material atomically coupled to the strip at the ends of the strip outside of a central sense region.

CROSS-REFERENCE TO RELATED APPLICATION This application claims priority under 35 U.S.C. §119(e) to Provisional Patent Application 60/840,123, filed Aug. 25, 2006, and titled “Digital Electronic Dispersion Compensation for Multi-Mode Fiber.” TECHNICAL FIELD This description relates to analog circuits. In particular, this description relates to an amplifier and associated circuit topology for achieving variable gain amplification with high bandwidth and fine granularity promoting high linearity. BACKGROUND In many applications it may be necessary to amplify an analog signal exhibiting a wide amplitude range. For example, a wide range of input signals may be present at the receiving end of a multi-mode fiber optic cable. Such a signal may require analog conditioning or digital signal processing to correct for degradation introduced by the physical medium of transmission, i.e., the optical cable itself. In many signal conditioning systems especially communication links, in order to compensate for a wide amplitude range of received information bearing signals, the input signals are subjected to amplitude adjustment using a VGA (“Variable Gain Amplifier”). A VGA allows for the selection and adjustment of gain to be applied to an input signal. Amplitude adjustment or so called gain adjustment of an incoming signal by a VGA is used to achieve an amplitude level well above the noise and offset thresholds. Without the application of gain adjustment, it may not be feasible to perform further post processing of an incoming signal, such as adaptive equalization. Cascading gain stages may provide a wide range of amplification and/or attenuation. However, each additional stage may be undesirable as it will introduce harmonic distortion. Harmonic distortion typically arises due to non-linearities inherent in each stage. Thus, it is desirable to devise an amplitude adjustment scheme using a VGA with a low number of gain stages such that the VGA is suitable for high bandwidth and high linearity applications with a wide amplitude adjustment range. SUMMARY According to one general aspect, a high bandwidth, a fine granularity variable gain amplifier (“VGA”) is described comprising a gain block, the gain block comprising at least one input node, a low gain output tap and a high gain output tap, a parallel gain block, wherein the parallel gain block comprises a low gain signal path and a high gain path, the low gain signal path and the high gain signal path respectively coupled to the low gain output tap and the high gain output tap of the attenuator, a cascaded gain block, wherein the low gain signal path and the high gain signal path are coupled to an input of the cascaded gain block, and a gain adjustment control to adjust a gain of the VGA, wherein the gain adjustment control is configured to cause a selective activation of at least a portion of the low gain path or the high gain path in the parallel gain path to achieve a desired overall gain. According to another general aspect, a method for providing variable gain amplification of an input signal with high bandwidth and high linearity is described comprising configuring a low gain signal path and a high gain signal path, receiving an input signal, passively generating a first attenuated signal and a second attenuated signal from the input signal, the first attenuated signal having a larger attenuation than the second attenuated signal, generating a low gain amplified signal from the first attenuated signal and a high gain amplified signal from the second attenuated signal via respectively the low gain signal path and the high gain signal path, generating a composite signal by combining the low gain signal and the high gain signal, and amplifying the composite signal to generate a VGA output signal. The details of one or more implementations are set forth in the accompanying drawings and the description below. Other features will be apparent from the description and drawings, and from the claims. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is a block diagram of a VGA for achieving a high gain output signal from a received input signal with high bandwidth and linearity. FIG. 2 is a flowchart illustrating example operations of the VGA topology of FIG. 1 . FIG. 3 is a schematic of a VGA with high bandwidth and high linearity that incorporates differential signaling. FIG. 4 is schematic of a cascaded gain stage. FIG. 5 is a schematic of a parallel gain block. FIG. 6 is a schematic of a differential pair that may be utilized as a gm element in a parallel gain stage or a cascade gain stage. FIG. 7 is a flowchart of a process for selecting a gain for a VGA. DETAILED DESCRIPTION FIG. 1 is a block diagram of a VGA topology 100 for achieving a high gain output signal 122 having substantially high bandwidth and linearity from a received input signal 120 . The VGA topology may receive an input signal 120 which is passed to a gain block 102 , a parallel gain block 112 , and a cascade gain block 110 to produce a VGA output signal 122 . Although FIG. 1 depicts single-ended signals, it will be understood by skilled practitioners that the topology shown in FIG. 1 may be utilized with differential signals. The gain block 102 may compensate for a wide range of input signal 120 amplitudes. The gain block 102 may provide gain greater than unity, in which case it may function as an amplifier. Alternatively, the gain block 102 may provide gain less than unity in which case it may function as an attenuator. According to one embodiment, the gain block 102 may be an attenuator that comprises passive components to achieve attenuation of the input signal with high bandwidth. For example, the gain block 102 may comprise a resistive ladder, including a plurality of resistors (described below with reference to FIG. 3 ). The gain block 102 may comprise a plurality of output taps (e.g., 124 , 126 ) that provide output signals at various attenuation amplitudes. For example, output tap 124 may provide a high gain output (ATT 1 ) feeding high gain signal path 120 while output tap 126 may provide a low gain output (ATT 2 ) feeding low gain signal path 122 . The parallel gain block 112 may include a plurality of parallel gain stages 104 , 106 . Although only two parallel gain stages 104 , 106 are shown in FIG. 1 , it will be understood by skilled practitioners that the parallel gain block may include an arbitrary number of parallel gain stages (e.g., 104 , 106 ). As described below, the parallel gain stages 104 , 106 may each respectively comprise a plurality of gm cells (not shown in FIG. 1 but described below). Gm refers to a transconductance of a simple amplification circuit in which a voltage signal is received at an input to generate a current signal at an output. Both parallel gain stages 104 , 106 may be identical comprising identical gm stages with identical current densities and to minimize phase delay between the respective outputs of the gain stages 104 , 106 . The parallel gain stages 104 , 106 in parallel gain block 112 may respectively be placed in the high gain signal path 120 and the low gain signal path 122 . In particular, the high gain signal path 120 may be coupled to the high gain output tap 124 of gain block 102 while the low gain signal path 122 may be coupled to the low gain output tap 126 of the gain block 102 . As described below, depending upon the particular combination of gm stages comprising each of the parallel gain stages 104 , 106 that may be selectively activated, various amplification levels may be achieved at summing block 108 . The number of combinations of gm stages that may be activated may directly provide fine granular control of the amplification level (e.g., each combination may provide a varying level of gain adjustment). A digital control block 124 may be utilized to control the activation of gm stages within parallel gain block 112 . One particular example of a process for selecting a gain for a VGA is illustrated below with respect to FIG. 7 . The outputs of the parallel gain block 112 (e.g., parallel gain stages 104 , 106 ) may be summed at a summing block 108 to provide an input to a cascade gain block 110 . Although FIG. 1 shows only two summed gain stages ( 104 , 106 ), it will be understood by skilled practitioners that any arbitrary number of parallel gain stages comprising a parallel gain block 112 may be summed at the summing block 108 . The summed output signals from the parallel gain block 112 may be received by a cascade gain block 110 where the summed signal is amplified by one or more cascaded gain stages (e.g., 114 ( 1 ), 114 ( 2 ), 114 ( 3 )). Although the cascade gain block 110 shown in FIG. 1 shows three cascaded gain stages 114 ( 1 ), 114 ( 2 ) and 114 ( 3 ), it will be understood by skilled practitioners that the cascade gain block 110 may include any number of cascaded gain stages. The output signal 122 of the cascade gain block 110 may then further processed. Inductive peaking may be utilized at the output of every gain stage in order to reduce phase delay problems and increase bandwidth. FIG. 2 is a flowchart illustrating example operations of the VGA topology of FIG. 1 . The process is initiated ( 202 ). A low gain signal path and a high gain signal path may be configured ( 204 ). As described with respect to FIG. 1 , the low gain signal path and high gain signal path may each respectively comprise a separate gain stage further, each gain stage further comprising a plurality of amplification elements. Each of the amplification elements may be selectively activated or deactivated via the digital control block 124 . An input signal that is to be amplified may be received ( 206 ). A low gain attenuated signal and a high gain attenuated signal may be respectively generated ( 208 ). The generation of the low and high gain attenuated signals may be achieved via the gain block 102 . A low gain amplified signal and a high gain amplified signal may be respectively generated from the low gain attenuated signal and the high gain attenuated signal via the low gain signal path and the high gain signal path ( 210 ). A composite signal may be generated by combining the low gain amplified signal and the high gain amplified signal ( 212 ). This combination may be achieved, for example, by summing the low gain amplified signal and the high gain amplified signal at a common node. The composite signal may then be further amplified to generate a VGA output signal ( 214 ). This further amplification may be achieved in a cascaded fashion. VGA amplification may now be completed ( 216 ). FIG. 3 is a schematic of a VGA amplifier 300 with high bandwidth and high linearity that incorporates differential signaling. FIG. 3 provides a specific example of the topology shown in FIG. 1 . The VGA may include a gain block 102 , a parallel gain block 112 and a cascade gain block 110 . A differential signal (not shown in FIG. 3 ) may be received via differential inputs Inp 302 and Inn 304 at gain block 102 . The gain block 102 may be an attenuator that comprises a resistive ladder having a plurality of resistors (e.g., R 1 304 , R 2 306 , R 3 308 , R 4 310 , R 5 312 and R 6 314 ). The gain block 102 may provide various output taps for generating differential output signals having various levels of attenuation. For example, the gain block 102 may generate high gain differential signals (ATT 1 P 316 , ATT 1 N 318 ) and low gain differential signals (ATT 2 P 322 , ATT 2 N 324 ). Thus, the gain block 102 may include output taps 316 generating signal ATT 1 P, output tap 318 generating signal ATT 1 N, output tap 322 generating signal ATT 2 P and output tap 324 generating signal ATT 2 N. In the case where the gain block 102 is an attenuator using passive components, the various amplitude output signals from the gain block 102 may be generated via the technique of a voltage divider using a resistive ladder as shown in FIG. 3 . The ATT 1 N and ATT 2 N signals shown in FIG. 3 may correspond to the following amplitudes: ATT ⁢ ⁢ 1 ⁢ P = R ⁢ ⁢ 2 + R ⁢ ⁢ 3 R ⁢ ⁢ 1 + R ⁢ ⁢ 2 + R ⁢ ⁢ 3 ⁢ ( Inp - Vcm ) ATT ⁢ ⁢ 2 ⁢ P = R ⁢ ⁢ 3 R ⁢ ⁢ 1 + R ⁢ ⁢ 2 + R ⁢ ⁢ 3 ⁢ ( Inp - Vcm ) Similarly the ATT 2 N and ATT 2 P signals shown in FIG. 3 may correspond to the following amplitudes: ATT ⁢ ⁢ 1 ⁢ N = R ⁢ ⁢ 2 + R ⁢ ⁢ 3 R ⁢ ⁢ 1 + R ⁢ ⁢ 2 + R ⁢ ⁢ 3 ⁢ ( Inn - Vcm ) ATT ⁢ ⁢ 2 ⁢ N = R ⁢ ⁢ 3 R ⁢ ⁢ 1 + R ⁢ ⁢ 2 + R ⁢ ⁢ 3 ⁢ ( Inn - Vcm ) High gain differential signal (ATT 1 P, ATT 1 N) and low gain differential signal (ATT 2 P, ATT 2 N) may be respectively provided to parallel block 112 . In particular, as shown in FIG. 3 the high gain differential signal (ATT 1 P, ATT 1 N) may be provided to a first differential parallel gain stage 326 while low gain differential signal (ATT 2 P, ATT 2 N) may be provided to a second differential parallel gain stage 328 . Each of the differential parallel gain stages ( 326 , 328 ) may provide amplification of their respective input signals. According to one embodiment, parallel gain stages 326 and 328 may be identical providing identical amplitude gain and phase delay. It is also possible to achieve a unity gain by omitting resistor R 1 . The differential outputs of the parallel gain block 112 may be provided to a cascade gain stage 110 where they are summed at respective common differential node inputs ( 338 , 340 ). The differential signal provided to cascade gain block 110 may then be amplified by any number of cascaded gain stages (e.g., 330 ( 1 ), 330 ( 2 ), 330 ( 3 ), 330 ( 4 )) to generate differential outputs 342 and 344 . Although FIG. 3 shows only four cascaded gain stages, it will be understood that cascade gain stage 110 may include any arbitrary number of cascaded gain stages. FIG. 4 is schematic of a cascaded gain stage 110 . Any number of cascaded gain stages 330 may be included in a cascade gain block 110 . Each cascaded gain stage 330 may include any number of gm elements 408 ( 1 )- 408 ( m ) (described below with reference to FIG. 6 ), which in aggregate operation achieve a gain for the cascaded gain stage 330 . In particular, each gm element 408 ( 1 )- 408 ( m ) may be selectively activated or deactivated to adjust an overall gain for the cascaded gain stage 330 . Each of the gm elements may be viewed as a separate amplification element for the stage. Each of the gm elements 408 ( 1 )- 408 ( m ) may be a differential pair arranged in a common source configuration as described below with reference to FIG. 6 . The gm elements 408 ( 1 )- 408 ( m ) may all respectively be coupled at their source nodes to a tail transistor 402 that may operate as a current source. Further, the gm elements 408 ( 1 )- 408 ( m ) may all respectively be coupled to a load block ZLOAD 416 , which may comprise either a passive load such as a resistor or an active load possibly generated using one or more MOSFET transistors. A voltage source AVDD 406 may be coupled to the load block ZLOAD 416 , to provide a voltage bias. The source of the tail transistor 404 may be coupled to a common voltage reference AVSS 404 . Differential input signal INP 1 338 and INN 1 340 may be provided as input to each of the gm elements 408 ( 1 )- 408 ( m ). The input differential signal INP 1 338 and INN 1 340 may be amplified by each of the gm elements 408 ( 1 )- 408 ( m ) to generate a composite amplified differential signal OUTP 1 410 and OUPTN 1 412 . FIG. 5 is a schematic of a parallel gain block 112 . The parallel gain block 112 may include a first differential parallel gain stage 326 and a second differential parallel gain stage 328 . Each parallel gain stage 326 , 328 may include any number of gm elements (respectively 408 ( 1 )- 408 (n) and 408 (n+1)- 408 (n+o)) (described below with reference to FIG. 6 ), which in aggregate operation achieve a gain for each of the parallel gain stages 326 and 328 in the parallel gain block 112 . In particular, each gm element 408 ( 1 )- 408 ( n ) and 408 ( n+ 1)- 408 ( n +o) may be selectively activated or deactivated to adjust an overall gain for their respective gain stages 326 , 328 . Each of the gm elements 408 ( 1 )- 408 ( n ) and 408 ( n+ 1)- 408 ( n+o ) may be a differential pair arranged in a common source configuration as described below with reference to FIG. 6 . The gm elements 408 ( 1 )- 408 ( n ), 408 ( n+ 1)- 408 ( n+o ) may all respectively be coupled to a tail transistor 514 that may operate as a current source. Further, the gm elements 408 ( 1 )- 408 ( n ), 408 ( n+ 1)- 408 ( n+o ) may all respectively be coupled to a load block ZLOAD 516 , which may comprise either a passive load such as a resistor or an active load possibly generated using one or more MOSFET transistors. A voltage source AVDD 406 may be coupled to the load block ZLOAD 516 , to provide a voltage bias. The source of the tail transistor 514 may be coupled to a common voltage reference AVSS 404 . Differential input signal INP 2 502 , INN 2 504 may be provided as input to the first parallel gain stage 326 , while differential input signal INP 3 518 , INN 3 520 may be provided as an input to the second differential parallel gain stage 328 . Each of the differential signals (INP 2 502 , INN 2 504 and INP 3 518 , INN 3 520 ) may be amplified by the respective parallel gain stages 326 , 328 to generate respective outputs (not shown in FIG. 5 ), which are combined at a common node to produce a differential output signal OUTP 2 510 , OUTN 2 512 . FIG. 6 is a schematic diagram of a differential pair that may be utilized as a gm element in a parallel gain stage or a cascade gain stage. The differential pair 600 may include a first input transistor 602 and a second input transistor 604 respectively receiving input signals Inp 606 and Inn 608 . Although the differential pair 600 shown in FIG. 6 utilizes nmos transistors for input transistors 602 , 604 , the input transistors 602 , 604 may also be pmos transistors. The substrate of each of the input transistors 602 , 604 may be coupled to a common substrate node Vsub 606 . FIG. 7 is a flowchart that illustrates a process for increasing or decreasing the gain achieved with a VGA. The process is initiated in step 702 . In step 704 , a set of gm elements in the low gain path 122 may be activated. In step 706 , the gain may be increased by turning on a set of gm elements in the cascade block 110 . In step 708 , the gain may be further increased by turning on gm elements in the high gain path 120 while correspondingly turning off the same number of gm elements in the low gain path 122 . By turning on and off corresponding numbers of gm elements in the low gain 122 and high gain paths 120 insures maintenance of a fixed common mode. Of course the overall gain for the VGA may be lowered by following the process as shown in FIG. 7 in reverse order. The source nodes of each of the input transistors 602 , 604 may both be coupled together at a common node 610 , which is also coupled to a third transistor Mena 612 . The third transistor Mena 612 may itself be coupled to a current source, for example a common current source transistor such as Mtail 514 shown in FIG. 5 . The differential pair 600 may be selectively activated or deactivated by respectively turning transistor Mena 612 on or off. Gain may be achieved for the differential pair 600 by steering current from the source/drain of transistor Inp 606 and the source/drain of transistor Inn 608 . The third transistor Mena 612 may receive bias signal Vena 614 , with its substrate node biased by a signal Vsub_ena 610 . The various techniques described herein may be implemented in digital electronic circuitry, or in computer hardware, firmware, software, or in combinations of them. Furthermore, these techniques may also be implemented as a computer program product, i.e., a computer program tangibly embodied in an information carrier, e.g., in a machine-readable storage device or in a propagated signal, for execution by, or to control the operation of, data processing apparatus, e.g., a programmable processor, a computer, or multiple computers. A computer program, such as the computer program(s) described above, can be written in any form of programming language, including compiled or interpreted languages, and can be deployed in any form, including as a stand-alone program or as a module, component, subroutine, or other unit suitable for use in a computing environment. A computer program can be deployed to be executed on one computer or on multiple computers at one site or distributed across multiple sites and interconnected by a communication network. Method steps may be performed by one or more programmable processors executing a computer program to perform functions by operating on input data and generating output. Method steps also may be performed by, and an apparatus may be implemented as, special purpose logic circuitry, e.g., an FPGA (field programmable gate array) or an ASIC (application-specific integrated circuit). Processors suitable for the execution of a computer program include, by way of example, both general and special purpose microprocessors, and any one or more processors of any kind of digital computer. Generally, a processor will receive instructions and data from a read-only memory or a random access memory or both. Elements of a computer may include at least one processor for executing instructions and one or more memory devices for storing instructions and data. Generally, a computer also may include, or be operatively coupled to receive data from or transfer data to, or both, one or more mass storage devices for storing data, e.g., magnetic, magneto-optical disks, or optical disks. Information carriers suitable for embodying computer program instructions and data include all forms of non-volatile memory, including by way of example semiconductor memory devices, e.g., EPROM, EEPROM, and flash memory devices; magnetic disks, e.g., internal hard disks or removable disks; magneto-optical disks; and CD-ROM and DVD-ROM disks. The processor and the memory may be supplemented by, or incorporated in special purpose logic circuitry. To provide for interaction with a user, implementations may be implemented on a computer having a display device, e.g., a cathode ray tube (CRT) or liquid crystal display (LCD) monitor, for displaying information to the user and a keyboard and a pointing device, e.g., a mouse or a trackball, by which the user can provide input to the computer. Other kinds of devices can be used to provide for interaction with a user as well; for example, feedback provided to the user can be any form of sensory feedback, e.g., visual feedback, auditory feedback, or tactile feedback; and input from the user can be received in any form, including acoustic, speech, or tactile input. Implementations may be implemented in a computing system that includes a back-end component, e.g., as a data server, or that includes a middleware component, e.g., an application server, or that includes a front-end component, e.g., a client computer having a graphical user interface or a Web browser through which a user can interact with an implementation, or any combination of such back-end, middleware, or front-end components. Components may be interconnected by any form or medium of digital data communication, e.g., a communication network. Examples of communication networks include a local area network (LAN) and a wide area network (WAN), e.g., the Internet. While certain features of the described implementations have been illustrated as described herein, many modifications, substitutions, changes and equivalents will now occur to those skilled in the art. It is, therefore, to be understood that the appended claims are intended to cover all such modifications and changes as fall within the true spirit of the various embodiments.

A fine granularity, wide-range variable gain amplifier (“VGA”) comprises an attenuator, a high gain signal path, a low gain signal path and a gain adjustment control to adjust a gain of the VGA, wherein the gain adjustment control is configured to cause a selective activation of at least a portion of the low gain signal path or the high gain signal path to achieve a desired overall gain.

BACKGROUND OF THE INVENTION [0001] 1. Field of the Invention: [0002] The present invention is related to secure storage of digital information. More specifically, the present invention is related to encryption and the encoding and decoding of digital information of, for example, an odometer. [0003] 2. Description of the Related Technology: [0004] With the advent of electronic odometers in vehicles, odometer fraud has become easier for educated criminals. Electric odometers store the odometer value in non-volatile memory, usually an EEPROM or similar device in order to ensure that the odometer value is preserved when the vehicle is switched off. The problem with the EEPROM is that it has a standard protocol interface that allows easy reading and writing to the EEPROM for someone educated in the electrical arts. [0005] A typical prior art electric odometer is illustrated in FIG. 1. The odometer 10 is composed of a microcontroller 12 that is used for processing the value for a display device (not shown). The serial EEPROM 14 is used to store the odometer value which can be written to the EEPROM or read from the EEPROM by the microprocessor 12 via signal lines 17 and 19 , respectively. As shown in FIG. 1, however, these signals could also be written/read via contacts 16 and 18 , respectively. [0006] There have been prior attempts to remedy the problems presented by electronic odometers. For example U.S. Pat. No. 5,924,057 disclosed a method of preventing odometer fraud employing a multiplex ring and a plurality of vehicle control modules. The odometer value was stored in the cluster control module and, periodically, in another module. If the main odometer value was tampered with, it would quickly be detected upon comparison with the value stored in an alternate module. [0007] U.S. Pat. No. 5,297,178 discloses a tamper resistant system having a programmable memory counter with a plurality of memory locations arranged sequentially. Each of the memory locations includes a predetermined number of storage elements. Each storage element is programmable from a first value to a second value but not visa-versa. As these storage elements comprise the counter, the odometer reading of a vehicle cannot be reduced. Neither of these solutions, however, completely protects the odometer value, which is still subject to tampering, albeit with greater difficulty. [0008] There is, therefore, a need in the art for a vehicle odometer that cannot be altered after initial setting. It is an object of the present invention to overcome the limitations inherent in the prior art. SUMMARY OF THE INVENTION [0009] The present invention solves the problems inherent in the prior art by providing a different vehicle odometer having an encoder in lieu of the standard EEPROM. [0010] The encoder utilizes an encryption algorithm. The present invention will work with any encryption algorithm, the selection of which is left for the customer. According to one embodiment, when the encoder receives an increment signal, such as one generated when a vehicle travels a pre-determined distance, an odometer value is combined with another value to form a larger (longer) value that is harder to break. The larger bit value may then be encrypted by the encoder using the encryption algorithm and stored in the encoder's own EEPROM. A microcontroller is used to receive the encrypted value from the encoder and then to decrypt it in real time for transmission to a display device. [0011] Other and further objects, features and advantages will be apparent from the following description of presently preferred embodiments of the invention, given for the purpose of disclosure and taken in conjunction with the accompanying drawings. BRIEF DESCRIPTION OF THE DRAWINGS [0012] [0012]FIG. 1 is a block diagram of a prior art electronic odometer; [0013] [0013]FIG. 2 is a block diagram of an electronic odometer of the present invention; [0014] [0014]FIG. 3 is a flowchart showing the method of the present invention; [0015] [0015]FIG. 4 is a flowchart of a sub-step of the method of the present invention; [0016] [0016]FIG. 5 is a flowchart of another sub-step of the method of the present invention; [0017] [0017]FIG. 6 is a block diagram illustrating a first alternate configuration of the present invention; [0018] [0018]FIG. 7 a is a block diagram illustrating a second alternate configuration of the present invention; [0019] [0019]FIG. 7 b is a block diagram illustrating a third alternate configuration of the present invention; [0020] [0020]FIG. 8 is a block diagram illustrating the functional components of an encoder of the present invention; and [0021] [0021]FIG. 9 is a flowchart further illustrating the method illustrated in FIG. 4. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS [0022] The basic configuration of the present invention is illustrated in FIG. 2. Quite simply, the electronic odometer 20 of the present invention is composed of a microcontroller 22 (having an integral decoder or an external decoder) and an encoder 24 that are connected for communication purposes by two signal lines 23 and 25 . While this description uses, as an illustrative example, a mileage odometer, it will be clear to those skilled in the relevant art that the present invention can be used for many applications that require a secure (encrypted) counter without departing from the scope and spirit of the appended claims (e.g., gas meters, electric meters, etc.). [0023] Alternate implementations of the encoder/decoder of the present invention are illustrated in FIGS. 6, 7 a, and 7 b. For example, FIG. 6 illustrates a serial configuration. According to FIG. 6, the encoder 64 receives an input signal from a voltage source that is responsive to a mileage sensor 69 . The mileage sensor 69 effectively opens and closes the switch 68 which, when closed, stimulates the encoder 64 to increment the mileage-related value. In this configuration, the encoder 64 communicates with the decoder 62 , which then decodes the encrypted signal to obtain the mileage value which the decoder 62 dutifully transmits to the display 66 as illustrated in FIG. 6. In an alternate embodiment of the present invention, the mileage sensor 69 , the switch 68 , and the encoder 64 can be integrated—making it more difficult to tamper with the odometer value stored within the encoder 64 . [0024] The mixed serial/parallel configuration illustrated in FIG. 7 a differs significantly from the configuration illustrated in FIG. 6. The encoder/decoder combination of FIG. 7 a most closely resembles the one found in FIG. 2. FIG. 7 a shows a mileage sensor 79 that sends an update signal (either analog or digital) to the decoder 72 . When the mileage sensor 79 generates a signal, the decoder 72 , if equipped sufficiently with circuitry and/or logic, can decide whether or not to update the encoder 74 . In an advantageous embodiment of the present invention, the internal configuration of the encoder is only 16 bits, which is insufficient resolution for typical automobile odometer applications. Hence, the advantageous embodiment uses an odometer-related value (known as the PEG value, which, for example, is updated every ten miles) in conjunction with a segment value that is incremented every mile (or other mileage sensor signal interval as the case may be). In this example, the following formula is used: actual odometer value={current PEG value×10}+segment value [0025] The configuration of FIG. 7 a is well suited to the 16-bit limitation of the preferred embodiment of the encoder 74 because the decoder 72 has sufficient circuitry/logic to update the encoder once every ten miles (or other peg period) rather than updating every mileage event as with the configuration of FIG. 6. In this illustration, the encoder 74 is considered to be the master by the odometer system. The configuration of FIG. 7 a also has an additional benefit over the simpler configuration of FIG. 6. The configurations of FIGS. 7 a and 7 b enables the decoder 72 and the encoder 74 to perform “mutual authentication” which is much more secure and much harder to temper because of the pure digital exchange of information that is preferred. [0026] [0026]FIG. 7 b is a modification of the configuration illustrated in FIG. 7 a. Security is enhanced by enhancing the functionality of the encoder and decoder. As illustrated in FIG. 7 b, encoder/processor/decoder (EPD) 73 is in functional communication with the corresponding decoder/processor/encoder (DPE) 75 . Although the configuration of FIG. 7 b functions in much the same manner as the configuration illustrated in FIG. 7 a, the additional circuitry/logic of EPD 73 and DPE 75 enable a complete encryption authentication scheme between both elements in both directions of information flow rather than the single direction (i.e., only from output from encoder 74 ) of FIG. 7 a. [0027] As mentioned previously, the odometer counter of the present invention is composed of an encoder 24 and a microcontroller 22 as shown in FIG. 2. An example of a suitable encoder 24 for purposes of the present invention is the HCS201 manufactured by Microchip Technology Inc. of Chandler, Ariz. The microcontroller 22 can be any processor or microcontroller that has sufficient memory and computing capacity to perform the decryption and/or data processing tasks. It should be noted from the outset that the present invention is not limited to a specific encryption algorithm. Any encryption algorithm will work with the present invention, the selection of which is left to the user. [0028] The encoder 24 may, according to the embodiment, be equipped with a unique key. While the key length is not terribly important to the function of the present invention, it is understood in the encryption art that the longer the key, the more secure the data protected by that key. In the preferred embodiment of the present invention, the key has at least 64 bits. Preferably, the key is used to configure an encryption engine. The encryption engine resides within the encoder 24 and can be composed of hardware or software or any combination there between. The encoder 24 of the present invention has sufficient computing capacity to execute the encryption scheme on the encryption engine using the odometer-related value to form an encrypted odometer-related value. According to one embodiment, the un-encrypted odometer-related value, once incremented, is then stored in a non-volatile memory in the encoder 24 . It should be noted that in the preferred embodiment of the present invention, the encoder 24 is unable to transmit anon-encrypted odometer related value externally. Thus, according to the preferred embodiment, a request for odometer call to the encoder 24 always returns an encrypted odometer value and thus dissuades tampering. Moreover, the encoder 24 is unable to accept a “write” function call that would overwrite the odometer-related data stored within a specific location of the non-volatile memory of the encoder 24 . Finally, as a necessary function required of an odometer, the encoder 24 is able to increment the odometer value upon receiving an increment signal from, for example, the microcontroller 22 . Once the odometer value is increased, it may be stored in the non-volatile memory of the encoder 24 . While various non-volatile memories can be used with the present invention, the advantageous embodiment of the present invention has the encoder 24 equipped with an EEPROM specific for this application with a single bit write. This non-volatile internal configuration aids greatly in error recovery, particularly when the present invention is used in harsh environments or when there is an unreliable power supply or frequent shut-downs. [0029] In addition to the encoder 24 , the odometer counter 20 of the present invention includes a microcontroller 22 . In the preferred embodiment of the present invention, the microcontroller 22 has sufficient computing capacity to decrypt the encrypted odometer-related value as received from the encoder 24 upon the proper read call being made to the encoder 24 , at which point the microcontroller 22 then processes the data. Finally, the microcontroller 22 is also equipped to transmit the decrypted odometer value to another device, such as a display device (not shown). [0030] It should be noted that the specific encryption algorithm used, be it implemented in hardware or software or some combination there between, is not directly important to the method of the present invention. It is envisioned that various users would chose different encryption algorithms, key lengths, and other implementation details to suit their particular needs. [0031] In the preferred embodiment of the present invention, the odometer value is only 16 bits long because the encoder is configured internally as a 16-bit device, although other configured devices can be used without departing from the spirit and scope of the present invention. This may create a problem in that 16 bits is not sufficient to contain a foreseeable mileage number. This problem, however, can be remedied, for example, by incrementing the counter only every ten miles and keeping track of the remaining miles in the microcontroller memory. If the power to the system is lost or broken, the maximum degradation of accuracy would be only ten miles. This degradation problem could be remedied, for example, by using a larger bit length for the odometer value. A 32-bit length (or more) would be far more than sufficient. However, a 32-bit length for the odometer value would also add considerable cost and complexity to the encoder 24 . An alternate (albeit more expensive) embodiment of the present invention utilizes a 32-bit encoder and would thus use, store, and encrypt a 32-bit long odometer value without the need for an additional 16-bit value as described in the preferred embodiment. An 8-bit or less configured device could be used; however, an 8-bit device would have a very low mileage limit, but may be useful in some applications. In short, alternate applications of the present invention may require longer bit configurations, while still other alternate applications could make do with shorter bit configurations. [0032] If the odometer value is only 16 bits in length, it may be necessary to add extra bits before the odometer value is encrypted. In the preferred embodiment of the present invention, a second 16 bit value is combined with a 16 bit odometer value to form a 32 bit concatenated value before the encryption process is executed. In the preferred embodiment of the present invention, the second 16 bit value is composed of a predetermined 12 bit value (that is retrieved from the encoder's nonvolatile memory) that is concatenated with a four bit value that is derived from inputs into the encoder 24 . There are several advantages to this configuration. First, this configuration enables the 16-bit value to be modified from counter system to counter system, thereby making each counter system 20 unique and thus harder to circumvent. Secondly, the first twelve bits of the sixteen bit “constant” can be a checksum (using, for example the CRC algorithm) combined with a scaling factor for the encoder counter. [0033] The scaling factor is the factor that determines when the odometer mileage PEG (e.g., the 10-mile increments) is incremented. These twelve bits also allow for the use of an embedded checksum within the encrypted transmission for post-transmission analysis (to ensure that there was no tampering of the encrypted data during transmission), as well as an incrementing factor. The incrementing factor is used to obtain the actual odometer mileage peg from the odometer-related value using the following formula: (Odometer-related value)*(Incrementing factor)=Odometer Mileage PEG [0034] In the preferred embodiment of the present invention, the non-encrypted odometer value is stored within a non-volatile memory of the encoder 24 (see FIG. 2). Thereafter, and for as many times as the microcontroller 22 requests, a serial number (identifying the specific encoder 24 ) is attached to the data packet containing the encrypted odometer value and then transmitted to the microcontroller 22 . In the preferred embodiment of the present invention, the transfer is asynchronous, although alternate transfer mechanisms are possible. The use of an identifier code, such as the serial number, precludes the substitution of the original encoder 24 with another encoder 24 (having a lower odometer value). An attacker can read the serial number but not the internal checksum data. Thus, the attacker cannot use the serial number to tamper with the odometer-related value without affecting the checksum. [0035] An overview of a method according to one embodiment of the present invention, is illustrated in FIG. 3. This method of the present invention presumes that there is an encoder having a unique key as well as encryption algorithm embodied either as hardware, software, or some combination therebetween. This sample method starts with step 302 . Initially (i.e., upon encoder startup), the unique key is used to configure the encryption engine, step 304 . Once initiated, the encoder is ready for operation. In typical operation, the encoder receives an increment signal, usually from the decoder, in step 306 . Next, in step 308 , the encoder increments the odometer value and then stores the incremented value back into the encoder's memory. After incrementing the odometer value, the encoder then encrypts the odometer value (or odometer-related value, as the case may be) in step 310 . Normally, the decoder updates a display device on a periodic basis (normally once every few seconds). To this end, the encoder, in step 312 , transmits the encrypted odometer value to the decoder, usually in an asymmetric fashion. Once the encrypted data is received by the decoder, the decoder first identifies the encoder from which it received the encrypted data to determine, through an identifying code with which each encoder is equipped, step 314 . Next, in step 316 , a check is made by the decoder to see whether or not the identifying code matches the one from the expected encoder. If not, the encoder is presumed to be a fake and execution skips to step 317 where an error message (indicative of tampering) is issued and the data discarded. Otherwise, the decoder decrypts the encrypted odometer value and sends the decrypted odometer value to the display device for display to the end user and the process ends in step 318 . In either case, the process ends in step 320 . [0036] Step 310 of a method according to one embodiment of the present invention, may be comprised of four sub-steps that are illustrated in FIG. 4. As mentioned previously, in the preferred embodiment, the odometer value is 16 bits in length. This 16-bit odometer-related value is first read from the encoder's protected memory in step 404 . Next, in step 406 , the odometer-related value is incremented and then the incremented odometer-related value is written back to the encoder's protected memory. In step 408 , the incremented odometer-related value is then added to another, according to this embodiment of the present invention, a 16-bit value to form a 32-bit packet of data. It should be noted that not all of the 16-bit value need be stored in the encoder's protected memory. Part of the 16-bit value can be stored elsewhere, such as in the decoder's non-volatile (protected) memory. Furthermore, part of the 16-bit value can be composed of a checksum of the odometer-related value combined with another value that is stored in the encoder's protected memory. Next, in step 410 , the 32-bit value is encrypted to form a 32-bit packet of encrypted data and the sub step ends in step 410 . [0037] Similarly, step 312 of FIG. 3 is, according to the preferred embodiment, itself composed of three sub steps that are illustrated in FIG. 5. The sub step begins in step 502 . Thereafter, a 32-bit data packet containing the 16 bit constant and the 16-bit odometer-related value is encrypted to generate an encrypted odometer-related value in step 504 . Next, in step 506 , the encrypted 32-bit data packet is combined with an identification code to form a tagged data packet. Finally, the tagged data packet is then transmitted from the encoder to the decoder in step 508 and the process ends in step 510 . [0038] [0038]FIG. 8 illustrates in block diagram form an embodiment of an encoder useful in the practice of the present invention. The encoder 800 has four input pins 802 , 804 , 806 , and 808 entitled S3, S2, S1, and S0, respectively. In addition to the four input pins, the encoder 800 also has a V SS pin 816 , an LED pin 812 , and V CC pin 810 . The output from the encoder 800 is through the data pin 814 . The encoder 800 is equipped with program/data memory and/or registers that are interconnected by signal lines and/or a common or separate busses. The actual internal configuration for transmitting data from program/data memory or non-volatile memory to other registers or memory, or the choice of signal lines and/or busses, is not important to the present invention. Any internal communication mechanism that can retrieve data from one location within the encoder and deliver it to another area of the microcontroller will suffice. [0039] The respective states of the four input pins 802 , 804 , 806 , and 808 are placed into a buffer 818 as illustrated in FIG. 8. The value from the buffer 818 forms the 4-bit function component 854 of the full data packet 850 . The value from the buffer 818 also forms a component of the 16-bit “constant” and is also stored in the 4-bit segment 826 of the pre-encryption buffer 820 . The encoder 800 is equipped with a non-volatile (protected) memory 830 . Within the non-volatile memory 830 are the seed variable 832 , the identifying code variable 834 , the encryption key 836 , some discrimination data 838 , and the counter 840 which contains the odometer-related value. The counter 840 can optionally be a synchronization counter. When the counter 840 needs to be incremented, the Gray coded counter 840 is read out of the memory 830 and is copied into the incrementer 842 where the Gray code is modified by one bit to an increase of one mileage segment. The incremented Gray code counter value 840 is then written back into the memory 830 . The other copy of the Gray coded counter 840 is sent to the Gray code-to-Sequential code converter 844 which converts the Gray coded counter value 840 into normal sequential code (i.e., standard binary digits) and then writes that sequential odometer-related value into a 16-bit segment 822 of the pre-encryption buffer 820 . The third component of the pre-encryption data is the discrimination data 838 constant that is copied from the non-volatile memory 830 to the discrimination data segment 824 of the pre-encryption buffer 820 . In alternate embodiments of the present invention, however, part of the constant data can be obtained from the memory location 838 and the remainder from some other place or, preferably, calculated in real-time, such as a CRC value based upon the current odometer related value stored in the 16-bit segment 822 . As mentioned before, the encryption key 836 can be stored into the memory 830 before first operation of the encoder 800 , or, it can be generated using the seed 832 singly or in conjunction with the discrimination data 838 . [0040] After all of the data segments 822 , 824 , and 826 have been loaded into the buffer 820 , the encryption key 836 is loaded into the encryption engine 828 as illustrated in FIG. 8. Within the encryption engine 828 , the 32 bits of clear text data are transformed into 32 bits of encrypted data which is then stored in the 32-bit encrypted data segment 858 of the full data packet 850 . Either before, during, or after the 32-bet encrypted value is loaded into the buffer segment 858 (the order does not matter), the other segments of the fall data packet 850 are appended (loaded) into their various segments 852 , 854 , and 856 . Status bits, (e.g., a low voltage bit and a repeat bit) which can indicate the functional condition of the encoder 800 , are optionally stored into the status segment 852 of the full data packet 850 . As mentioned before, the contents of the buffer 818 containing the state of the four input bits may, optionally, be loaded into the 4-bit segment 854 . It is important to note that adding the four input bits may pose a security hazard, as it would give an opportunity for an attacker to manipulate the four input signals (and thus the ultimate output data signal) in order to modify the odometer value. Thus, the overall system would be more secure if the four input values were not added to the full data packet, however remote the possibility of manipulation. [0041] In addition to the other values, the identification code 834 is read out of the non-volatile memory 830 and copied into the serial segment 856 of the full data packet 850 . The complete 66-bit full data packet is then transmitted out through the data pin 814 to the decoder. [0042] [0042]FIG. 9 illustrates a detailed description of the method of the present invention. The method of FIG. 9 starts at step 902 . The process is initialized (started) in step 904 . Thereafter, a request to process is detected by the encoder in step 906 . Next, in step 908 , the counter value is read from the non-volatile protected memory. Normally, the counter value (i.e., the odometer-related value) is in the form of Gray code. A binary Gray code with n digits corresponds to a Hamiltonian path on an n-D hypercube (including direction reversals). For the present invention, a binary Gray code is a number encoding scheme such that adjacent numbers have a single bit difference. Use of the Gray code with the present invention simplifies the processing and the implementation scheme. [0043] In step 912 , the incremented counter (the incremented odometer-related value) is first stored in a buffer and a copy is written back to the encoder's non-volatile protected memory. Next, in step 914 , a 12-bit constant or portion thereof is read from the encoder's protected memory. A 4-bit function code that is comprised of the status of each of the four input pins is then attached to the 12-bit value to form a 16-bit value in step 916 . Next, the 16-bit value is catenated with the 16-bit odometer-related value to form a 32-bit value that is stored into a buffer. Thereafter, according to one embodiment of the present invention, the encryption engine is configured with a unique encryption key, thereby forming a unique encryption engine, step 920 . The 32-bits of plain data are then encrypted to obtain 32-bits of encrypted data forming a 32-bit data packet, step 922 . Typically, this encryption process uses a one-to-one function (to maintain the same number of bits) although other functions can be used without departing from the scope of the present invention. Next, a 28 bit identifying code, such as a serial number, is attached to the 32-bit. Optionally, a status code, and/or the function code mentioned previously can be attached to the identification code to form a fall data packet, step 924 . In step 926 , the full data packet is transmitted to the decoder. [0044] According to one embodiment, upon receipt of the full data packet from the encoder, the decoder identifies the source of the full data packet using either the identifying code or a function of the identifying code in step 928 , and, if found to originate from the correct encoder, decrypts the encrypted data packet and validates the odometer-related value with the decryption checksum in step 930 . Finally, the PEG value, which is first multiplied by the scaling factor, is then added to the segment value in order to obtain the actual odometer value, which is then used to update any electronic and/or display modules in step 932 . The process then ends at step 934 . [0045] The present invention, therefore, is well adapted to carry out the objects and attain both the ends and the advantages mentioned, as well as other benefits inherent therein. While the present invention has been depicted, described, and is defined by reference to particular advantageous embodiments of the invention, such references do not imply a limitation on the invention, and no such limitation is to be inferred. The invention is capable of considerable modification, alternation, alteration, and equivalents in form and/or function, as will occur to those of ordinary skill in the pertinent arts. The depicted and described advantageous embodiments of the invention are exemplary only, and are not exhaustive of the scope of the invention. Consequently, the invention is intended to be limited only by the spirit and scope of the appended claims, giving full cognizance to equivalents in all respects.

An improved vehicle odometer is provided with an encoder and a microcontroller. The encoder has a unique key that configures an encryption algorithm. The encryption algorithm can be in the form of circuitry or software or any combination there between. The encoder receives a increment signal from the vehicle indicating that the vehicle has traveled a pre-determined distance. The encoder then increments the odometer-related value and then encrypts it with the encryption algorithm. The encoder then wraps the encrypted odometer-related value into a data packet with a serial number. Only a suitably equipped computing device, such as a microcontroller, identify and correctly decrypt and process the encrypted odometer-related value, thereby preventing the tampering of the vehicle's odometer.

BACKGROUND OF THE INVENTION [0001] The present invention relates to a gate turn-off thyristor that use a wide-gap semiconductor and relates, in particular, to a gate turn-off thyristor capable of interrupting a large current within a wide temperature range. [0002] As a first background art gate turn-off thyristor (hereinafter referred to as GTO) that uses silicon, there is the one disclosed in JP H06-151823 A. In the first background art GTO, a mesa-type p-type base layer is provided on an n-type base layer that has an anode electrode, and an n-type emitter layer is formed by impurity diffusion in a central region of the mesa-type p-type base layer. With this construction, a junction between the p-type base layer and the n-type emitter layer is not exposed on the mesa slope, and therefore, a GTO in which electric field concentration hardly occurs on the mesa slope can be obtained. However, since the n-type emitter layer is formed by impurity diffusion, there are many crystal defects, and the on-resistance of the GTO is increased. [0003] A second background art GTO that uses silicon is disclosed in JP 2692366 B. In the second background art, an n-type base layer is formed on a p-type emitter layer, and a p-type base layer is formed on the n-type base layer. An n-type emitter layer is formed by impurity diffusion on the p-type base layer, and a mesa-type n-type emitter layer is obtained by etching. The second background art is the same as the first background art regarding the point that the n-type emitter layer is formed by impurity diffusion. [0004] As a third background art gate turn-off thyristor (hereinafter abbreviated to GTO) that uses a wide-gap semiconductor such as silicon carbide (SiC), there is, for example, the one described on pages 518 through 520 of the reference document: IEEE Electron Device Letters, Vol. 18, No. 11, November, 1997. In this background art, a p-type anode emitter layer is etched into a mesa-type down to a p-type base layer with which the anode emitter layer is put in contact, and a gate electrode is provided on the base layer so as to surround the anode emitter layer etched into the mesa-type. The structure is presumably adopted for the reasons as follows. In a GTO of silicon (Si) that is not the wide-gap semiconductor, a partial pn junction is generally formed by impurity thermal diffusion or ion implantation. However, in the case of SiC that is the wide-gap semiconductor, the impurity thermal diffusion is very slow and therefore not appropriate for mass production. Therefore, the pn junction is formed by ion implantation. In the case, if high-concentration impurity ions are implanted, the crystal defects increase and the resistance becomes high. Therefore, when a large current is flowed through the GTO, a voltage drop in the region where ions have been implanted increases, and the power loss is large. In particular, when impurity ions of a large atomic radius of a p-type impurity of aluminum or the like are implanted, crystal defects easily occur, and a high-concentration p-type region cannot be formed without a crystal defect. Accordingly, when a partial pn junction is formed at an SiC and particularly when a p-type region that flows a large current is formed, a p-type epitaxial film that has a good crystallinity and a little crystal defect is formed on an n-type layer. A GTO is formed by selectively etching the epitaxial film and forming a mesa-type partial pn junction. An end portion of the junction between the p-type layer and the n-type layer is exposed in the neighborhood of a mesa slope or a mesa corner portion. By covering the entire surface of the GTO after film formation with an insulator, ions from the outside are prevented from adhering to the semiconductor surface, and the long-term reliability of the GTO is secured. [0005] In general, the GTO has a current controllability to effect turn-off by diverting the principal current into the gate by applying a reverse bias voltage between the gate and the anode. Characteristics that represent the controllability include a “maximum controllable current”. The maximum controllable current means the maximum current that the GTO can control. In order to increase the maximum controllable current of the GTO, the principal current is diverted into the gate as much as possible by raising the off-gate voltage (reverse voltage applied between the gate and the anode) at the turn-off time. It is known that the maximum controllable current can be increased as the principal current to be diverted into the gate is increased by raising the off-gate voltage. [0006] FIGS. 13 and 14 show sectional views of typical GTO's of the second and third background art using SiC, respectively. In the second background art GTO shown in FIG. 13 , a lightly doped p-type SiC base layer 2 is formed on a heavily doped n-type SiC cathode emitter layer 1 that has a cathode electrode 21 connected to a cathode terminal K (hereinafter referred to as a cathode K) on its lower surface. An n-type base layer 3 is formed on the p-type base layer 2 . A p-type layer, a central region of which is left and becomes a p-type anode emitter layer 4 in a subsequent process, is formed by the epitaxial growth method on the entire surface of the n-type base layer 3 . Next, the mesa-type anode emitter layer 4 is formed by etching away a region of the p-type layer other than a region that becomes the anode emitter layer 4 by the reactive ion etching method until the surface of the n-type base layer 3 is somewhat removed. An n-type gate contact region 6 is formed by ion implantation so as to surround the anode emitter layer 4 in a portion located apart from a junction J of the end portion of the exposed n-type base layer 3 . An anode electrode 20 connected to an anode terminal A (hereinafter referred to as an anode A) is formed on the anode emitter layer 4 , and a gate electrode 22 connected to a gate terminal G (hereinafter referred to as a gate G) is formed on the gate contact region 6 . Finally, in order to prevent moisture and ions of Na ions and the like from adhering to the surface of the GTO, an insulator 10 of silicon dioxide (SiO 2 ) or the like is formed on the entire surface excluding the electrodes. [0007] The third background art GTO shown in FIG. 14 has substantially the same construction as that of the GTO shown in FIG. 13 except that the conductive types of the layers and the regions are inverted from those of the GTO shown in FIG. 13 . [0008] In the GTO shown in FIG. 13 , an off-gate voltage is applied between the gate G and the anode A at the turn-off time. Moreover, in the GTO shown in FIG. 14 , an off-gate voltage is applied between the cathode K and the gate G at the turn-off time. As a result, the principal current is diverted into the gate G to turn off the GTO both in the GTO's of FIGS. 13 and 14 . [0009] If the off-gate voltage is raised in order to increase the maximum controllable current in the GTO of FIG. 13 , an electric field at the insulator 10 in the neighborhood of the end region T of the junction J between the anode emitter layer 4 and the n-type base layer 3 is increased. In the case of SiC, the dielectric breakdown field is about ten times that of Si, and therefore, the thickness of the base layer 3 is reduced to several tenths of the thickness of Si. Therefore, if the off-gate voltage is raised, a high electric field is applied to the insulator 10 (e.g., SiO 2 film) on the surface of the mesa that forms the anode emitter layer 4 , and this might cause the dielectric breakdown of the insulator 10 . Moreover, there is a problem that, if the high electric field is continuously applied for a long term, a leakage current increases to reduce the gate withstand voltage (withstand voltage between the gate G and the anode A) of the GTO element, and the long-term reliability is degraded. [0010] Also, in the GTO of FIG. 14 , if the off-gate voltage is raised as a countermeasure for increasing the maximum controllable current, the electric field at the insulator 10 in the neighborhood of the end region T of the junction J between the cathode emitter layer 24 and the base layer 5 is increased. Consequently, the withstand voltage between the cathode K and the gate G is lowered, and the long-term reliability is degraded. [0011] As another countermeasure for increasing the maximum controllable current, a method for reducing the resistance in the transverse direction by increasing the impurity concentration of the base layer on which the gate electrode is provided and a method for increasing the thickness of the base layer are described in JP S61-182260 A. If the resistance in the transverse direction of the base layer is reduced by increasing the impurity concentration, the implantation efficiency of carriers (e.g., holes in the case of the GTO of FIG. 13 or electrons in the case of the GTO of FIG. 14 ) implanted from the emitter layer located adjacent to the base layer at the time of turning on the GTO is reduced. Moreover, if the base layer on which the gate is provided is increased in thickness, the amount of carriers, which move from the adjacent emitter layer through the base layer to the underlying base layer, is reduced. As a result, a gate current necessary for turning on the GTO is increased. Moreover, the on-state voltage is also raised, which causes a problem that the power loss is increased. [0012] The maximum junction temperature during the use of a semiconductor device that uses a wide-gap semiconductor is significantly higher than the maximum junction temperature (about 125° C.) during the use of a semiconductor that uses an Si semiconductor. For example, the maximum junction temperature during the use of SiC is not lower than 500° C. Therefore, in a device that uses a wide-gap semiconductor, the semiconductor device should preferably maintain the desired characteristics within a wide temperature range from room temperature to a temperature of not lower than 500° C. [0013] According to a background art reference of Material Science Forum Vols. 389-393 (2002), pp. 1349-1352, it is disclosed that the maximum controllable current is significantly reduced when the use temperature becomes 150° C. or higher in the GTO of SiC. For example, at a temperature of 200° C., the maximum controllable current is reduced to about one sixth or less of the maximum controllable current at room temperature. This is presumably for the reasons as follows. [0014] For the sake of easy understanding, the case of the GTO of Si is described first. In the case of Si, boron or aluminum is used as an acceptor. The substances have shallow acceptor levels of 45 meV and 60 meV and are easily ionized at room temperature, generating holes from the acceptor. Therefore, almost all the impurities are ionized, generating holes. When Si is used at a temperature up to the maximum junction temperature of 125° C., the ionization rate scarcely poses a problem since the impurity ionization rate is sufficiently high. [0015] Boron and aluminum, which are also used as an acceptor in the case of GTO of SiC as in the case of GTO of Si, have deep acceptor levels of about 300 meV and about 240 mV, respectively, and a low ionization rate of not higher than several percent at room temperature. However, the ionization rate is significantly increased when the temperature is elevated. [0016] For example, in the GTO of FIG. 13 , when the temperature is elevated to 150° C. or higher and the ionization rate of the p-type anode emitter layer 4 is increased, the number of holes implanted into the p-type base layer 2 from the anode emitter layer 4 via the n-type base layer 3 becomes significantly greater than at room temperature. Moreover, since electrons are also increased and the excess carriers (holes and electrons) are increased in the p-type base layer 2 , the maximum controllable current is reduced. Furthermore, since the carrier lifetime also becomes longer at a high temperature of not lower than 150° C., the maximum controllable current is also significantly reduced by this. Moreover, since the carrier density in the p-type anode emitter layer 4 is increased at high temperature, the depletion layer does not sufficiently spread when the off-gate voltage is applied. In such a state, the electric field is increased in the neighborhood of the end region T of the anode emitter layer 4 in the neighborhood of the junction J between the p-type anode emitter layer 4 and the n-type base layer 3 , and the withstand voltage (about 30 V) between the anode A and the gate G is lowered. [0017] Moreover, if the GTO of FIG. 14 enters a state as described above, the withstand voltage (about 30 V) between the cathode K and the gate G is lowered. Furthermore, the electric field in the neighborhood of the end region T of the cathode emitter layer 24 is increased, and the electric field of the insulator 10 is increased, possibly causing dielectric breakdown. Moreover, the leakage current is increased, and this reduces the reliability during long-term use. SUMMARY OF THE INVENTION [0018] According to the present invention, in the gate turn-off thyristor (hereinafter referred to as a wide-gap GTO) of a wide-gap semiconductor that has a mesa-type emitter layer, the maximum controllable current is increased by relieving the electric field of the insulator located in the neighborhood of the end portion of the junction between an emitter layer and a base layer where a gate is provided adjacent to the emitter layer. [0019] In order to relieve the electric field of the insulator located in the neighborhood of the end portion of the junction, a low-resistance gate region of a low resistance value is formed in the base layer. With this arrangement, a current at the turn-off time flows through the low-resistance gate region of a low resistance value, and therefore, a voltage drop is a little. If the gate current is increased at the turn-off time by raising the off-gate voltage, the electric field of the insulator is not increased so much. As another method for preventing the increase in the electric field of the insulator, there is a method for forming a field relief region in the neighborhood of the junction. Since the electric field of the insulator is relieved by the method, the off-gate voltage can be raised. Therefore, the principal current can be diverted with high efficiency. Since the off-gate voltage can be raised, a large maximum controllable current can be maintained within a wide temperature range from a low temperature of not higher than room temperature to an elevated temperature that exceeds 500° C. When the off-gate voltage is not increased so much, the long-term reliability is remarkably improved. Since the electric field of the insulator in the neighborhood of the junction can be reduced, the long-term reliability of the GTO can be maintained. [0020] A gate turn-off thyristor of a wide-gap semiconductor of the present invention includes a first emitter layer of either one of n-type and p-type conductive types having a first electrode on its one surface and a first base layer of a conductive type different from that of the first emitter layer provided on the other surface of the first emitter layer. This gate turn-off thyristor further includes a second base layer of a conductive type identical to that of the first emitter layer provided on the first base layer, a mesa-type second emitter layer of a conductive type different from that of the first emitter layer provided on the second base layer and a second electrode provided on the mesa-type second emitter layer. A low-resistance gate region is provided so as to surround the mesa-type second emitter layer in a region located apart from an end portion of a junction between the mesa-type second emitter layer and the second base layer, formed in a region that extends from a neighborhood of the end portion of the junction to a bottom portion of the mesa-type second emitter layer with interposition of the second base layer between the region and the junction, and having a conductive type identical to that of the second base layer and an impurity concentration higher than that of the second base layer. A third electrode is put in contact with an end portion of the low-resistance gate region. [0021] According to the present invention, by virtue of the first conductive type low-resistance gate region formed in the first conductive type base layer, an electron current flows from the first conductive type base layer through the first conductive type low-resistance gate region and the first conductive type gate contact region to the gate at the turn-off time. Since the low-resistance gate region has a low resistance value, a voltage drop in the first conductive type base layer is small even when a current due to the electron flow is large. Therefore, the off-gate voltage applied between the anode and the gate is not influenced by the voltage drop, and a large current can be flowed with high efficiency. As a result, the controllable current of the GTO can be increased. [0022] In another aspect of the present invention, a gate turn-off thyristor of a wide-gap semiconductor includes a first emitter layer of either one of n-type and p-type conductive types having a first electrode on its one surface and a first base layer of a conductive type different from that of the first emitter layer provided on the other surface of the first emitter layer. This GTO further includes a second base layer of a conductive type identical to that of the first emitter layer provided on the first base layer and a mesa-type second emitter layer of a conductive type different from that of the first emitter layer provided on the second base layer. This GTO further includes a second electrode, which is put in contact with a central region of the mesa-type second emitter layer and put in contact with the second emitter layer via a contact electrode in a region excluding the central region of the second emitter layer. In a region located apart from the end portion of the junction between the mesa-type second emitter layer and the second base layer, a low-resistance region, which has a conductive type identical to that of the second base layer and an impurity concentration higher than that of the second base layer, is provided so as to surround the mesa-type second emitter layer. A third electrode is provided in contact with the end portion of the low-resistance region. [0023] In another aspect of the present invention, a gate turn-off thyristor of a wide-gap semiconductor includes a first emitter layer of either one of n-type and p-type conductive types having a first electrode on its one surface, a first base layer of a conductive type different from that of the first emitter layer provided on the other surface of the first emitter layer, a second base layer of a conductive type identical to that of the first emitter layer provided on the first base layer, a mesa-type second emitter layer of a conductive type different from that of the first emitter layer provided on the second base layer, a high-resistance region provided in a central region of an upper surface of the second emitter layer and having a conductive type identical to that of the second emitter layer and an impurity concentration lower than that of the second emitter layer, a contact electrode put in contact with the second emitter layer and the high-resistance region, a second electrode put in contact with the contact electrode and the second emitter layer at least at a peripheral portion of the contact electrode, and having a contact resistance to the emitter layer greater than a contact resistance of the contact electrode to the emitter layer, a low-resistance region provided so as to surround the mesa-type second emitter layer in a region located apart from an end portion of a junction between the mesa-type second emitter layer and the second base layer, and having a conductive type identical to that of the second base layer and an impurity concentration higher than that of the second base layer, and a third electrode put in contact with an end portion of the low-resistance region. [0024] According to the present invention, the second electrode is put in contact with only the central region of the second emitter layer and put in contact with the second emitter layer via the contact electrode formed of a material of which the contact resistance to the semiconductor layer is lower than that of the second electrode in the other region. Therefore, the contact resistance between the second electrode and the second emitter layer in the region located with interposition of the contact electrode is lower than that of the other region. With this arrangement, a current that flows from the second electrode into the second emitter layer flows more in the peripheral region located with interposition of the contact electrode than in the central region where the resistance is relatively high. In the GTO, the current control effect by virtue of the low-resistance region is great in a portion located near the low-resistance region, but the effect is reduced in the central region remote from the low-resistance region. In the present invention, the greater part of the electrification current is flowed in the peripheral region where the current control effect by virtue of the low-resistance region is high, so that the current in the central region of a low control effect is reduced. As a result, the efficiency of extracting a current from the gate at the turn-off time is increased, and therefore, the controllable current of the GTO is increased. [0025] According to the present invention, by sufficiently separating the gate contact region of the GTO that uses a wide-gap semiconductor from the junction between the mesa-type emitter layer and the base layer, the electric field in the neighborhood of the junction or in the neighborhood of the mesa corner portion is not increased even when the off-gate voltage is high. By raising the off-gate voltage, the current flowing between the anode and the cathode can efficiently be diverted into the gate, and the controllable current of the GTO can be increased. Moreover, since a high electric field is not applied to the insulator, the leakage current is not increased, and the long-term reliability can be maintained. [0026] By forming the low-resistance region adjacent to the gate contact region, a voltage drop caused by the current that flows through the low-resistance region at the turn-off time can be reduced. Therefore, even when the off-gate voltage is the same as the conventional one, the turn-off current can be diverted into the gate with higher accuracy than in the conventional GTO. Even if the p-type impurity ionization rate is increased than at room temperature or the carrier lifetime become longer during use at high temperature, the off-gate voltage can be raised. Furthermore, the gate current at the turn-off time can be diverted into the gate with high efficiency by the low-resistance region. Therefore, a GTO, which has a large controllable current within a wide temperature range from a low temperature of not higher than room temperature to a temperature that exceeds 500° C. and is able to maintain high reliability for a long term, can be provided. BRIEF DESCRIPTION OF THE DRAWINGS [0027] FIG. 1 is a top view of a gate turn-off thyristor of a first embodiment of the present invention; [0028] FIG. 2 is a sectional view of the gate turn-off thyristor of the first embodiment of the present invention; [0029] FIG. 3 is a sectional view of the gate turn-off thyristor of the second embodiment of the present invention; [0030] FIG. 4 is a sectional view of the gate turn-off thyristor of the third embodiment of the present invention; [0031] FIG. 5 is a sectional view of the gate turn-off thyristor of the fourth embodiment of the present invention; [0032] FIG. 6 is a sectional view of the gate turn-off thyristor of the fifth embodiment of the present invention; [0033] FIG. 7 is a sectional view of the gate turn-off thyristor of the sixth embodiment of the present invention; [0034] FIG. 8 is a sectional view of the gate turn-off thyristor of the seventh embodiment of the present invention; [0035] FIG. 9 is a sectional view of the gate turn-off thyristor of the eighth embodiment of the present invention; [0036] FIG. 10 is a sectional view of the gate turn-off thyristor of the ninth embodiment of the present invention; [0037] FIG. 11 is a sectional view of the gate turn-off thyristor of the tenth embodiment of the present invention; [0038] FIG. 12 is a sectional view of the gate turn-off thyristor of the eleventh embodiment of the present invention; [0039] FIG. 13 is a sectional view of a background art gate turn-off thyristor; and [0040] FIG. 14 is a sectional view of another background art gate turn-off thyristor. DETAILED DESCRIPTION OF THE INVENTION [0041] Preferred embodiments of the gate turn-off thyristor (hereinafter, abbreviated to GTO) that uses silicon carbide (hereinafter, SiC) of the present invention will be described with reference to FIGS. 1 through 12 . FIG. 1 is a top view of one element of the GTO of the first embodiment of the present invention. FIGS. 2 through 8 are sectional views of one element (unit) of the GTO of each embodiment of the present invention. No hatching is shown in the cross sections of the figures for the sake of clear views of the figures. The top views of the GTO's of the embodiments are basically similar to the one shown in FIG. 1 . In the actual construction of the GTO of each of the embodiments, a number (normally several tens to several thousands) of elements are connected together on an identical substrate horizontally in the figure, and the anode electrodes, gate electrodes and cathode electrodes are connected parallel or in series depending on the situation. Although not shown in the figures, it is general that a number of elements are horizontally arranged and a number of elements are also arranged in columns in the vertical direction in the figures. First Embodiment [0042] A GTO that uses SiC of the first embodiment of the present invention is described with reference to FIGS. 1 and 2 . FIG. 1 is a top view that shows the upper surface before the provision of an insulator 10 of the GTO of the first embodiment. FIG. 2 is a sectional view taken along the line II-II of FIG. 1 . In FIGS. 1 and 2 , the GTO of the present embodiment has a heavily doped cathode emitter layer 1 (first emitter layer) of an n-type (first conductive type) SiC semiconductor that has a thickness of about 350 μm and an impurity concentration of not smaller than about 10 19 cm −3 and is provided with a cathode electrode 21 (first electrode) connected to the cathode terminal K (cathode K, hereinafter). A lightly doped base layer 2 (first base layer) of a p-type (second conductive type) SiC semiconductor that has a thickness of 50 μm and an impurity concentration of about 10 16 to 10 13 cm −3 is formed on the cathode emitter layer 1 . A thin n-type base layer 3 (second base layer) of a thickness of several micrometers is formed on the p-type base layer 2 . A p-type layer that serves as a p-type anode emitter layer 4 is grown by the epitaxial growth method on the entire surface of the n-type base layer 3 leaving a central region in a subsequent process. Next, a mesa-type anode emitter layer 4 is formed by deeply etching the other region by the reactive ion etching method to an extent that the surface of the n-type base layer 3 is exposed and its surface portion is somewhat removed leaving a region that becomes the p-type anode emitter layer 4 (second emitter layer). By implanting ions into the exposed surface n-type base layer 3 , an n-type low-resistance gate region 5 and an n-type gate contact region 6 are successively formed so as to surround the anode emitter layer 4 . The impurity concentration of the low-resistance gate region 5 should preferably be three times the impurity concentration of the base layer 3 . The low-resistance gate region 5 may be formed down to the neighborhood of the upper surface of the base layer 2 in the ion implantation process. The low-resistance gate region 5 is formed slightly apart from the junction J between the anode emitter layer 4 and the base layer 3 . The gate contact region 6 is a low-resistance region of an impurity concentration higher than that of the low-resistance gate region 5 and is formed in a position located far apart from the junction J. An anode electrode 20 (second electrode) connected to an anode terminal A (anode A, hereinafter) is formed on the anode emitter layer 4 , and a gate electrode 22 (third electrode) connected to a gate terminal G (gate G, hereinafter) is formed on the gate contact region 6 . Finally, in order to prevent moisture and ions of Na ions and the like from adhering to the surface of the GTO after the formation of the layers, an insulator 10 of silicon dioxide (SiO 2 ) or the like is formed on the entire surface excluding the electrodes. Nitrogen can be used as an n-type impurity. Moreover, boron and aluminum can be used as a p-type impurity. [0043] The structural feature of the GTO of the present embodiment resides in that the n-type gate contact region 6 in the n-type base layer 3 is separated from the junction J between the p-type anode emitter layer 4 and the n-type base layer 3 , providing a large creeping distance. Another feature resides in that the low-resistance gate region 5 that has a low resistance value and a high n-type impurity concentration is formed in the n-type base layer 3 and separated a prescribed distance apart from the junction J between the p-type anode emitter layer 4 and the n-type base layer 3 in the direction of the base layer 2 . A distance between the n-type gate contact region 6 and an end portion JE of the junction J between the p-type anode emitter layer 4 exposed on the slope of the mesa M and the n-type base layer 3 is about 2 μm to 10 μm. Moreover, the n-type low-resistance gate region 5 is formed at a depth of about 0.3 μm to 5 μm from the upper surface of the n-type base layer 3 . [0044] In the present embodiment, as shown in FIG. 2 , the n-type low-resistance gate region 5 should preferably be formed only below the bottom surface MB of the mesa M. However, it is acceptable to extend a little the n-type low-resistance gate region 5 to the inside of the n-type base layer 3 below the p-type anode emitter layer 4 as indicated by an extension SE in FIG. 2 . Extending the low-resistance gate region 5 toward the extension SE increases the maximum controllable current, the minimum firing current and the on-state voltage of the GTO and lowers the withstand voltage. Conversely, making the low-resistance gate region 5 shorter than the length shown in FIG. 2 reduces the maximum controllable current, the minimum firing current and the on-state voltage of the GTO and raises the withstand voltage. Particularly, when a high withstand voltage is required, it is preferable to shorten the n-type low-resistance gate region 5 and separate the region from the anode emitter layer 4 . [0045] Operation of the GTO of the present embodiment is described below. If a forward bias voltage is applied between the anode A and the gate G by making the voltage of the gate G lower than the voltage of the anode A in a state in which the voltage of the anode A is higher than the voltage of the cathode K, then a current flows from the anode A to the gate G. In this state, holes are injected from the anode emitter layer 4 into the n-type base layer 3 and enter the p-type base layer 2 , while electrons are injected from the n-type cathode emitter layer 1 into the p-type base layer 2 , consequently turning on the GTO and putting it into the on-state. If a reverse bias voltage is applied between the anode A and the gate G, an electron current that flows from the cathode K to the anode A is diverted into the gate G, then the GTO is turned off. [0046] In the GTO of the present embodiment, the n-type gate contact region 6 is located apart from the junction J between the mesa-type p-type anode emitter layer 4 and the n-type base layer 3 . Therefore, the withstand voltage between the gate G and the anode A is not determined by the creeping distance between the two but determined by the dielectric breakdown field originally possessed by SiC in the p-type anode emitter layer 4 . Since SiC has a high dielectric breakdown field, the GTO of the present embodiment has a high withstand voltage. Moreover, even when the ionization rate of the p-type anode emitter layer 4 is increased and the hole concentration is increased at an elevated temperature, the high withstand voltage can be maintained. [0047] By virtue of the formation of the n-type low-resistance gate region 5 in the n-type base layer 3 , an electron current flows from the n-type base layer 3 to the gate G through the n-type low-resistance gate region 5 and the n-type gate contact region 6 at the turn-off time. Since the n-type low-resistance gate region 5 has a high impurity concentration and a low resistance value, the voltage drop in the n-type base layer 3 is small and the electric field applied to the insulator 10 in the neighborhood of the junction J is not increased even when the electron current is large. Therefore, the off-gate voltage applied between the anode A and the gate G is not influenced so much by the voltage drop, and the off-gate voltage can be raised. By raising the off-gate voltage, a large electron current can be flowed with high efficiency. As a result, the controllable current of the GTO of the present embodiment can be increased. When the GTO of the present embodiment is used at a high temperature of about 500° C., the maximum controllable current can be increased by raising the off-gate voltage in the GTO of the present embodiment even when the amount of holes injected into the n-type base layer 3 is increased as a consequence of an increase in the hole concentration of the p-type anode emitter layer 4 or when the lifetime of holes and electrons become longer due to the temperature rise. Since the electric field applied to the insulator in the neighborhood of the junction J can be reduced, long-term reliability can be maintained. [0048] In a concrete example of the GTO of the present embodiment, the withstand voltage between the gate G and the anode A was 150 V, which means that a remarkable rise in the withstand voltage was able to be achieved in comparison with the voltage of about 30 V of the GTO of the background art examples shown in FIGS. 9 and 10 . No high electrical field is applied to the insulator 10 in the neighborhood of the junction J even when the gate voltage at the turn-off time is raised to increase the maximum controllable current, and dielectric breakdown is hard to occur. Since no high electrical field is applied, a decrease of the withstand voltage between the gate G and the anode A induced by an increase in the leakage current between the gate G and the anode A does not arise, and high reliability can be maintained for a long term. [0049] In the present embodiment, the anode emitter layer 4 is formed by the epitaxial growth method. Since crystal defects are very little by the epitaxial growth method, holes can sufficient be injected into the n-type base layer 3 . Therefore, the on-state voltage is reduced to a low voltage of 3.7 V, and the loss can be reduced. For example, when an anode emitter layer was formed by the ion implantation method causing many crystal defects, the on-state voltage was 7.5 V. [0050] Although an angle of the slope of the mesa-type anode emitter layer 4 with respect to the surface of the base layer 3 is about 105 degrees in the example shown in FIG. 2 , the present invention is applicable when the angle is within a range of about 140 degrees to 55 degrees. Second Embodiment [0051] FIG. 3 is a sectional view of a GTO that uses SiC of the second embodiment of the present invention. In FIG. 3 , the p-type and the n-type of the layers are interchanged in the GTO of the present embodiment in comparison with the GTO of the first embodiment shown in FIG. 2 . A lightly doped n-type SiC base layer 2 (second base layer) that has a thickness of about 50 μm is formed on the upper surface of a p-type anode emitter layer 4 A (first emitter layer) that has a thickness of about 350 μm and is provided with an anode electrode 20 (first electrode) connected to the anode A on its lower surface. A thin p-type base layer 3 A (second base layer) that has a thickness of several micrometers is formed on the base layer 2 A, and an n-type layer of which central region is left in a subsequent process to serve as an n-type cathode emitter layer 1 A is formed by the epitaxial growth method on the entire surface of the p-type base layer 3 A. Next, a region is deeply etched by the reactive ion etching method to an extent that the surface of the p-type base layer 3 A is exposed and its surface portion is somewhat removed leaving the other region that becomes the cathode emitter layer 1 A (second emitter layer) of the n-type layer, forming the mesa-type cathode emitter layer 1 A. Then, the cathode electrode 21 (second electrode) is formed on the cathode emitter layer 1 A. A low-resistance gate region 5 A that has a p-type high impurity concentration by ion implantation and a low resistance and a p-type gate contact region 6 A are formed successively layered on the exposed p-type base layer 3 A so as to surround the cathode emitter layer 1 A. A gate electrode 22 (third electrode) is formed on the gate contact region 6 A. Finally, an SiO 2 insulator 10 is formed on the entire surface excluding the electrodes. [0052] In the GTO of the present embodiment, the gate electrode 22 and the cathode electrode 21 are adjacently located. Therefore, if a forward bias voltage is applied between the cathode K and the gate G in a state in which the voltage of the anode A is higher than the voltage of the cathode K, then a current flows from the gate G to the cathode K. As a result, holes are injected from the anode emitter layer 4 A into the n-type base layer 2 A and enters the p-type base layer 3 , while electrons are injected from the n-type cathode emitter layer 1 A into the p-type base layer 3 A, by which the GTO is turned on and put into the on-state. If a reverse bias voltage is applied between the cathode K and the gate G to divert the current that flows from the anode A to the cathode K into the gate G, then the GTO is turned off. [0053] In the GTO that uses SiC of the present embodiment, by virtue of the formation of the n-type low-resistance gate region 5 A in the p-type base layer 3 A, the current that flows from the anode A to the gate G at the turn-off time passes through the low-resistance gate region 5 A and the gate contact region 6 A. Since the low-resistance gate region 5 A has a low resistance value, a voltage drop is small, and a large current can be flowed through the gate G. Therefore, by operation substantially similar to that of the first embodiment, the electric field of the insulator 10 in the neighborhood of the junction J between the n-type cathode emitter layer 1 A and the p-type base layer 3 A can be reduced at the turn-off time and in the off-state. Moreover, by raising the off-gate voltage at the turn-off time, almost the same maximum controllable current as that at room temperature can be obtained even at an elevated temperature. Third Embodiment [0054] FIG. 4 is a sectional view of a GTO that uses SiC of the third embodiment of the present invention. In the GTO of the present embodiment shown in the figure, a p-type region 7 , which includes at least the neighborhood of the end portion of the junction J between the p-type anode emitter layer 4 and the n-type base layer 3 and expands from the neighborhood of a corner portion MC of the mesa M toward the gate electrode 22 , is formed in the n-type base layer 3 . The other construction is the same as that of the GTO of the first embodiment shown in FIG. 2 . By virtue of the formation of the p-type region 7 , the field intensity of the insulator 10 in the neighborhood of the mesa corner portion MC located at the end portion of the junction J between the p-type anode emitter layer 4 and the n-type base layer 3 can be relieved even when the off-gate voltage at the turn-off time is increased. As a result, the withstand voltage between the gate G and the anode A can be raised, and the controllable current can be increased. Moreover, since the intensity of the electric field applied to the insulator 10 can be reduced, the deterioration of the insulator 10 can be prevented. Therefore, an increase in the leakage current between the gate G and the anode A is not caused even during a long-term use, and high reliability can be maintained for a long term. In the concrete example of the GTO of the present embodiment, the withstand voltage between the gate G and the anode A was 205 V, which means that a withstand voltage higher than the withstand voltage (150 V) of the GTO of the first embodiment was able to be obtained. Fourth Embodiment [0055] FIG. 5 is a sectional view of a GTO that uses SiC of the fourth embodiment of the present invention. In the GTO of the present embodiment shown in the figure, the n-type low-resistance gate region 5 is provided in a portion of an end region of the n-type base layer 3 excluding the p-type anode emitter layer 4 . The n-type low-resistance gate region 5 is formed by self-alignment in the n-type base layer 3 by means of a mask for mesa etching for forming the p-type anode emitter layer 4 . Therefore, a process for forming the pattern of the n-type low-resistance gate region 5 can be eliminated. In the present embodiment, as in the third embodiment, a p-type region 7 , which includes at least the neighborhood of the end portion of the junction J between the p-type anode emitter layer 4 and the n-type base layer 3 and expands from the neighborhood of the corner portion MC of the mesa M to the gate electrode 22 , is formed in the n-type low-resistance gate region 5 . The other construction is the same as that of the first embodiment shown in FIG. 2 . The formation of the p-type region 7 prevents the formation of a junction between the heavily doped p-type anode emitter layer 4 and the heavily doped n-type low-resistance gate region 5 in the neighborhood of the mesa corner portion MC and forms a junction at the bottom surface of the mesa M. As a result, the field intensity of the insulator 10 in the neighborhood of the mesa corner portion MC is relieved, and the off-gate voltage can be raised. It is also acceptable to enlarge the p-type region 7 so that the region covers the mesa corner portion MC and form the region connected to the anode emitter layer 4 . According to the concrete example of the present embodiment, the withstand voltage between the gate G and the anode A was 130 V, which means that a withstand voltage remarkably higher than the withstand voltage (about 30 V) of the conventional GTO was able to be obtained although the withstand voltage is lower than the withstand voltage (150 V) of the GTO of the first embodiment. Since the withstand voltage is high, the controllable current can be increased by raising the gate voltage. In the GTO of the present fourth embodiment, the gate contact region 6 A, which is provided in the GTO of each of the first through third embodiments, is not provided. Therefore, the construction is simple, and the manufacturing cost is low. Fifth Embodiment [0056] FIG. 6 is a sectional view of a GTO that uses SiC of the fifth embodiment of the present invention. In FIG. 6 , the GTO of the present embodiment has a construction in which the n-type low-resistance gate region 5 is excluded from the GTO of the third embodiment shown in FIG. 4 . In a process for forming the n-type low-resistance gate region 5 in the GTO of the third embodiment, a heavily doped n-type layer is formed by carrying out ion implantation into the n-type base layer 3 . At this time, crystal defects easily occur in the n-type base layer 3 and the n-type low-resistance gate region 5 . As a result, a leakage current on a surface between the gate and the anode is increased. Since the n-type low-resistance gate region 5 is not provided in the present embodiment, the problem caused by the crystal defects does not occur in the n-type base layer 3 . [0057] As in the GTO of the third embodiment shown in FIG. 4 , in the present embodiment, a p-type region 7 , which includes at least the neighborhood of the end portion of the junction J between the p-type anode emitter layer 4 and the n-type base layer 3 and expands from the neighborhood of the corner portion MC of the mesa M toward the gate electrode 22 for relieving the field intensity, is formed in the n-type base layer 3 . As in the case of the fourth embodiment, also in the present embodiment, it is acceptable to enlarge the p-type region 7 so that the region covers the mesa corner portion MC and is connected to the anode emitter layer 4 . With this construction, the field intensity in the neighborhood of the mesa corner portion MC is not increased even if the off-gate voltage is raised, and therefore, the field intensity of the insulator 10 is not increased. Therefore, the deterioration of the insulator 10 is avoided. In a concrete example of the GTO of the present embodiment, the withstand voltage between the gate electrode 22 and the anode electrode 21 was 210 V. Since the off-gate voltage can be raised, a GTO of a large controllable current can be provided. Sixth Embodiment [0058] FIG. 7 is a sectional view of a GTO that uses SiC of the sixth embodiment of the present invention. In the GTO of the present embodiment, a p-type base layer 2 is formed by the epitaxial growth method on a heavily doped n-type SiC cathode emitter layer 1 provided with a cathode electrode 21 on its lower surface. Next, an n-type low-resistance gate region 5 is formed in both end regions of the p-type base layer 2 . Next, an n-type base layer and a p-type anode emitter layer, which become the n-type base layer 3 and the p-type anode emitter layer 4 , respectively, through a subsequent process, are successively layered by the epitaxial growth method on the entire surfaces of the p-type base layer 2 and the n-type low-resistance gate region 5 , respectively. Both the end regions of the n-type base layer and the p-type anode emitter layer are etched by the reactive ion etching method until the surface of the n-type low-resistance gate region 5 is exposed, forming the mesa-type n-type base layer 3 and p-type anode emitter layer 4 that have a mesa slope MS. According to the construction of the present embodiment, a junction JE exposed on the mesa slope MS of the p-type anode emitter layer 4 and the n-type base layer 3 is separated from the neighborhood of the mesa corner portion MC where electric field concentration easily occurs, allowing a sufficient creeping distance to be provided. Therefore, since the off-gate voltage can be raised, a GTO of a large controllable current can be provided. When the n-type low-resistance gate region 5 is formed in the n-type base layer 3 by implanting ions to a deep portion in the base layer 3 as in the cases of the first through fourth embodiments, crystal defects easily occur in the n-type base layer 3 . In contrast to this, when the n-type base layer 3 is formed by the reactive etching method, no crystal defect occurs in the n-type base layer 3 . Since the ion implantation is carried out only in forming the n-type low-resistance gate region 5 in the p-type base layer 2 , ion implantation processes are few, and the manufacturing processes of the GTO can be simplified. [0059] In the case of the GTO of SiC, the impurity concentration of the n-type base layer 3 is higher than in, for example, the GTO of Si. Therefore, a depletion layer does not spread so much in the n-type base layer 3 at the turn-off time. Therefore, the electric field concentration on the end portion of the low-resistance gate region, which causes a problem in the GTO of Si, does not occur. Therefore, the withstand voltage between the anode electrode 20 and the cathode electrode 21 can be raised. Seventh Embodiment [0060] FIG. 8 is a sectional view of a GTO that uses SiC of the seventh embodiment of the present invention. In the GTO of the present embodiment, at least one n-type low-resistance gate small region 55 is formed in the neighborhood of the surface of the p-type base layer 2 within an active region where a principal current flows between the n-type low-resistance gate regions 5 located at both end portions. The other construction is the same as that shown in FIG. 7 . In the GTO of the present embodiment, the greater part of the flow of electrons injected from the n-type cathode emitter layer 1 into the p-type base layer 2 at the turn-off time can effectively be diverted into the right and left gate electrodes 22 by the n-type low-resistance gate region 55 formed in the active region. With this arrangement, a GTO of a large controllable current can be provided. Even if the lifetime of the carriers (electrons and holes) becomes long at high temperature when the use temperature exceeds 150° C. or when the amount of holes that pass through the n-type base layer 3 and flows into the p-type base layer 2 are increased as a consequence of an increase in the hole density due to an increase in the ionization rate of the p-type anode emitter layer 4 , the controllable current is scarcely reduced. [0061] The present invention can also be applied to GTO's that are constituted by interchanging the n-type layers and regions with p-type layers and regions and interchanging the p-type layers and regions with n-type layers and regions in the first through seventh embodiments. Eighth Embodiment [0062] FIG. 9 is a sectional view of a GTO that uses SiC of the eighth embodiment of the present invention. In the figure, an anode contact electrode 61 is formed in a region excluding the center portion of the upper surface of the mesa-type anode emitter layer 4 . An anode electrode 60 connected to the anode A is put in contact with only the central region of the anode emitter layer 4 . In a peripheral region excluding the central region of the anode emitter layer 4 , the anode electrode 60 is put in contact with the anode emitter layer 4 via the anode contact electrode 61 . Nickel is used for the anode contact electrode 61 , and gold, of which the contact resistance to the semiconductor layer is higher than nickel, is used for the anode electrode 60 . In the present technical field, it is known that, when a metal film is formed on an SiC semiconductor layer, a contact resistance between the two is varied depending on the kind of the metal and heat treatment after the film formation of the metal film besides the electrical conductivity of the metal. Metals of a low contact resistance include nickel, titanium, aluminum, tungsten and composite films of these metals. Metals of a high contact resistance include gold and so on. In the present embodiment, the anode contact electrode 61 of nickel is provided divided into at least right and left two regions and subjected to appropriate heat treatment. A gap between the anode contact electrodes 61 located in the two regions is about 1 μm to 20 μm. The other construction is similar to that of the fifth embodiment shown in FIG. 6 . Although an angle between the side surface of the mesa-type anode emitter layer 4 and the surface of the base layer 3 is about 90 degrees in FIG. 9 , the present embodiment is also applicable even when the angle is within a range of about 140 degrees to 50 degrees. [0063] In the present embodiment, by virtue of a low contact resistance used for the anode contact electrode 61 , the contact resistance between the anode contact layer and the anode emitter layer 4 is reduced. Therefore, an on-state current (Hall current) scarcely flows through the region where the anode electrode 60 is put in direct contact with the anode emitter layer 4 and flows to the anode emitter layer 4 through the portion of the anode contact electrode 61 . Therefore, the current flows intensively under the anode contact electrode 61 , and a current density in the portion where the anode contact electrode 61 is not located is reduced. [0064] The electron current, which flows from the cathode emitter layer 1 , flows through the region where the Hall current is flowing, and therefore, the electron current also leans to the region where the anode contact electrode 61 is located. Therefore, a region of the electron current, where electrons exist in surplus, comes close to the gate contact region 6 . Therefore, electrons can efficiently be extracted from the gate G at the turn-off time, and the controllable current is increased. In the case of the present embodiment, the controllable current was increased by 55% in comparison with that of the standard background art example. [0065] Although the present embodiment has had the construction in which the anode contact electrode 61 is provided divided, the anode electrode 60 may be similarly divided. Moreover, a similar effect can be obtained even when the anode electrode 60 is not provided and only the anode contact electrode 61 is provided so long as no problem occurs in terms of bonding. Ninth Embodiment [0066] FIG. 10 is a sectional view of a GTO that uses SiC of the ninth embodiment of the present invention. In the figure, an anode contact electrode 71 of a material of a low contact resistance such as nickel is formed on the anode emitter layer 4 . A lightly doped p-type region 73 is provided in the central region of the surface of the anode emitter layer 4 . An anode electrode 70 is formed on the anode contact electrode 71 . The anode electrode 70 is made larger than the anode contact electrode 71 , and its end portion is put in direct contact with a peripheral portion of the anode emitter layer 4 . The other construction is similar to that of the eighth embodiment. Since the p-type region 73 has a low impurity concentration, a contact resistance between the region 73 and the anode contact electrode 71 is high. Moreover, the internal resistance of the region 73 is also higher than the peripheral anode emitter layer 4 . As a result, Hall current flows while being diverted to the right and left going around the center portion, as in the eighth embodiment. Therefore, the electron current also flows while being diverted into the right and left. The region, where electrons exist in surplus due to the diverted electron flows, comes close to the gate contact region 6 . Therefore, the electron density is reduced in the central region of the p-type base layer 2 . As a result, electrons can efficiently be extracted from the gate, and therefore, the controllable current is increased. Although the angle between the side surface of the mesa-type anode emitter layer 4 and the surface of the base layer 3 is about 90 degrees in FIG. 10 , the present embodiment is also applicable even when the angle is within a range of about 140 degrees to 50 degrees. Tenth Embodiment [0067] FIG. 11 is a sectional view of a GTO that uses SiC of the tenth embodiment of the present invention. The GTO of the present embodiment differs from that of the ninth embodiment in that a heavily doped n-type region 83 is formed in place of the lightly doped p-type region 73 . The other construction is similar to that of the ninth embodiment shown in FIG. 10 . Also, in the construction, a current scarcely flows in the central region of the anode emitter layer 4 , and Hall current flows while being diverted into the right and left in the anode emitter layer 4 . Therefore, an electron current also flows while being diverted into the right and left regions. As a result, the paths of the diverted electron flows come close to the gate contact region 6 . Therefore, the control efficiency is improved, and the controllable current is improved. The n-type region 83 can be concurrently formed by ion implantation when the gate contact region 6 is formed, and therefore, the formation processes become simplified. Eleventh Embodiment [0068] FIG. 12 is a sectional view of a GTO that uses SiC of the eleventh embodiment of the present invention. In the present embodiment, a heavily doped n-type region 93 is provided in a central region of the surface of the n-type base layer 3 . The other construction is similar to that of the ninth embodiment shown in FIG. 10 . In the construction, the flow of Hall current, which flows from the anode emitter layer 4 toward the cathode emitter layer 1 , is disturbed by the n-type region 93 formed in the central region of the base layer 3 and diverted into the right and left in the base layer 3 . As a result, the controllable current can be increased by an effect similar to that of the eighth embodiment. INDUSTRIAL APPLICABILITY [0069] The present invention can be used for the gate turn-off thyristor that uses a wide-gap semiconductor capable of interrupting a large current within a wide temperature range.

A mesa-type wide-gap semiconductor gate turn-off thyristor has a low gate withstand voltage and a large leakage current. Since the ionization rate of P-type impurities greatly increases at high temperatures when compared with that at room temperature, the hole implantation amount increases and the minority carrier lifetime becomes longer. Consequently, the maximum controllable current is significantly lowered when compared with that at room temperature. To solve these problems, a p-type base layer is formed on an n-type SiC cathode emitter layer which has a cathode electrode on one surface, and a thin n-type base layer is formed on the p-type base layer. A mesa-shaped p-type anode emitter layer is formed in the central region of the n-type base layer. An n-type gate contact region is formed sufficiently apart from the junction between the p-type anode emitter layer and the n-type base layer, and an n-type low-resistance gate region is so formed in the n-type base layer that it surrounds the anode emitter layer.

CROSS-REFERENCE TO RELATED APPLICATIONS [0001] This application is a continuation of U.S. application Ser. No. 12/552,972, filed Sep. 2, 2009, which is a divisional of U.S. application Ser. No. 12/129,574, filed May 29, 2008, which is a divisional of U.S. application Ser. No. 10/389,721 filed on Mar. 14, 2003, now U.S. Pat. No. 7,381,210, the contents of which are expressly incorporated herein by reference in their entirety. BACKGROUND OF THE INVENTION [0002] In vertebrate animals, the heart is a hollow muscular organ having four pumping chambers: the left atrium, the left ventricle, the right atrium and the right ventricle. The atria are isolated from their respective ventricles by one-way valves located at the respective atrial-ventricular junctions. These valves are identified as the mitral (or bicuspid) valve on the left side of the heart, and tricuspid valve on the right side of the heart. The exit valves from the left and right ventricles are identified as the aortic and pulmonary valves, respectively. [0003] The valves of the heart are positioned in valvular annuluses that comprise dense fibrous rings attached either directly or indirectly to the atrial and ventricular muscle fibers. Valve leaflets comprising flexible collagenous structures are attached to, and extend inwardly from, the annuluses to meet at coapting edges. The aortic, tricuspid and pulmonary valves each have three leaflets, while the mitral valve only has two. In normal operation, the leaflets of the mitral valve open as left ventricle dilates thereby permitting blood to flow from the left atrium into the left ventricle. The leaflets then coapt (i.e. close) during the contraction cycle of the left ventricle, thereby preventing the blood from returning to the left atrium and forcing the blood to exit the left ventricle through the aortic valve. Similarly, the tricuspid valve regulates flow from the right atrium into the right ventricle, and the pulmonary valve regulates blood exiting the right ventricle. [0004] For a number of clinical reasons various problems with heart valves can develop. One common form of heart disease involves the deterioration or degradation of the heart valves which leads to stenosis and/or insufficiency. Heart valve stenosis is a condition in which the valve does not open properly. Insufficiency is a condition in which the valve does not close properly. Insufficiency of the mitral valve, most common because of the relatively high fluid pressures in the left ventricle, results in mitral valve regurgitation (“MR”), a condition in which blood reverses its intended course and flows “backward” from the left ventricle to the left atrium during ventricular contractions. [0005] A number of surgical techniques have been developed to repair degraded or otherwise incompetent heart valves. A common procedure involves replacement of a native aortic or mitral valve with a prosthetic heart valve. These procedures require the surgeon to gain access to the heart through the patient's chest (or possibly percutaneously), surgically remove the incompetent native heart valve and associated tissue, remodel the surrounding valve annulus, and secure a replacement valve in the remodeled annulus. While such procedures can be very effective, there are shortcomings associated with such replacement valves. For example, the invasive nature of the implantation procedure typically results in substantial patient discomfort and requires patients to remain hospitalized for extended recovery periods. In addition, the two basic types of commercially available replacement valves, mechanical valves and tissue valves, each have shortcomings of their own. Mechanical replacement valves typically offer extended operational lifetimes, but the patient is usually required to maintain a regimen of anti-coagulant drugs for the remainder of his or her life. Tissue valves typically offer a higher degree of acceptance by the body which reduces or eliminates the need for anti-coagulants. However, the operational lifetimes of tissue valves is typically shorter than mechanical valves and thus may require a subsequent replacement(s) during the patient's lifetime. [0006] As an alternative to prosthetic heart valve replacement, it is often preferable to remodel the native heart valve and/or the surrounding tissue. Remodeling of the valve often preserves left ventricular function better than mitral valve replacement because the subvalvular papillary muscles and chordae tendineae are preserved (most prosthetic valves do not utilize these muscles). Valvular remodeling can be accomplished by implanting a prosthetic ring (a.k.a. “annuloplasty ring”) into the valve annulus to reduce and/or stabilize the structure of the annulus in order to correct valvular insufficiency. Annuloplasty rings are typically constructed of a resilient core covered with a fabric sewing material. Annuloplasty procedures can be performed alone, or they can be performed in conjunction with other procedures such as leaflet repair. Although such annuloplasty procedures have become popular and well accepted, reshaping the surrounding annulus and traditional leaflet repairs do not always lead to optimum leaflet coaptation. As a result, some patients may still experience residual mitral valve regurgitation following such annuloplasty procedures. [0007] A recently developed technique known as a “bow-tie” repair has also been advocated for repairing insufficient heart valves, in particular the mitral valve. The mitral valve bow-tie technique involves suturing the anterior and posterior leaflets together near the middle of their coapting edges, thereby causing blood to flow through two newly formed openings. While this does reduce the volume of blood that can flow from the atrium to the ventricle, this loss is compensated by improved leaflet coaptation which reduces mitral regurgitation. This process as originally developed by Dr. Ottavio Alfieri involved arresting the heart and placing the patient on extracorporeal bypass and required invasive surgery to access and suture the leaflets together. More recently, however, some have advocated a “beating heart” procedure in which the heart is accessed remotely and remains active throughout the bow-tie procedure. [0008] One particular method for performing a beating heart bow-tie procedure (i.e. without extracorporeal bypass) has been proposed by Dr. Mehmet Oz, of Columbia University. (See PCT publication WO 99/00059, published Jan. 7, 1999, the contents of which are incorporated herein by reference). In one embodiment of this procedure, the associated device consists of a forceps-like grasper used to grasp and hold the mitral valve leaflets in a coapted position prior to the connecting step. Since the mitral valve leaflets curve toward and slightly into the left ventricular cavity at their mating edges, the grasper device is passed through a sealed aperture in the apex of the left ventricle. The edges of the mating mitral valve leaflets are then grasped and held together, and subsequently a fastening device such as a clip or suture is utilized to fasten them. The Mehmet Oz disclosure also discloses teeth on the grasper device that are linearly slidable with respect to one another so as to permit alignment of the mitral valve leaflets prior to fastening. Since the procedure is done on a beating heart, it will be readily understood that the pressures and motions within the left ventricle and mitral valve leaflets are severe and render Dr. Oz's procedure very skill-intensive. [0009] The bow-tie technique has proved to be a viable alternative for treating otherwise incompetent heart valves. Nonetheless, shortcomings associated with the current bow-tie procedures have been identified. Current systems typically include tissue stabilizing devices having mechanical graspers, barbed members, and vacuum devices. Often, use of these devices results in the less than optimal leaflet stabilization and fastener placement. Many of these problems arise from the fact that the surgeon is required to capture, retain and fasten the leaflets in one relatively inflexible procedure. These difficulties are compounded when the leaflets are small or calcified making them difficult to pull together, and in beating heart procedures in which the leaflets are actively functioning. In addition, the size and complexity of most current devices make minimally invasive surgical procedures more difficult, if not impossible. In light of the foregoing, there is presently a need for improved systems for stabilizing multiple tissue heart valve leaflets and placing a fastening device therebetween. More specifically, there is a present need for an improved bow-tie procedure for repairing a patient's mitral valve. BRIEF SUMMARY OF THE INVENTION [0010] The present invention solves the problem of effectively stabilizing at least one tissue portion in vivo. Additionally, the present invention provides a device capable of delivering a fastener to the stabilized tissue portion through a catheter from a remote insertion location. [0011] In one aspect, the present invention is directed to a system for repairing tissue within the heart of a patient and includes a guide catheter having a proximal end, a distal end, and at least one internal lumen formed therein, a therapy catheter capable of applying at least one suture to the tissue, and a fastener catheter capable of attaching at least one fastener to the suture. The therapy catheter and the fastener catheter are capable of traversing the internal lumen of the guide catheter. [0012] In another aspect, the present invention pertains to a system for repairing tissue within the heart of a patient and comprises a guide catheter having a proximal end, a distal end, and at least one internal lumen formed therein, a therapy catheter having at least one needle lumen in communication with at least one needle port positioned therein, at least one needle positioned within the needle lumen, and a fastener catheter having at least one fastener detachably coupled thereto. In addition, the fastener catheter includes at least one cutting member. [0013] In yet another aspect, the present invention discloses a system for repairing tissue within the heart of a patient and includes a guide wire capable of being inserted into the patient and advanced through a circulatory pathway, a therapy catheter attachable to the guide wire and capable of applying at least one suture to the tissue, and a fastener catheter attachable to the guide wire and capable of attaching at least one fastener to the suture. [0014] In a further aspect, the present invention pertains to a guide catheter for delivering a tissue repair device to tissue located within the heart of a patient and comprises an outer wall defining an outer wall lumen, a directing lumen capable of receiving a steering device therein and a flexible support device positioned within the outer wall lumen. [0015] In another aspect, the present invention discloses a catheter for delivering a suture to tissue within the heart of a patient and includes an elongated body having a distal end, at least one suction recess formed on the distal end, at least one needle port located proximate to the suction recess, at least one needle lumen having at least one needle positioned therein in communication with the needle port, at least one needle receiving port having at least one needle catch located therein positioned proximate to the suction recess, and at least one actuator member in communication with the needle. [0016] In yet another aspect, the present invention is directed to a catheter for delivering a suture to tissue within the heart of a patient and comprises an elongated body having a distal end with at least one suction recess formed thereon, at least one needle port located proximate to the suction recess, at least one needle lumen having at least one detachable needle attached to suture material positioned therein and in communication with the needle port, at least one needle receiving port located proximate to the suction recess, at least one needle trap capable of receiving the detachable needle positioned within the needle receiving port, and at least one actuator member in communication with the needle. [0017] In yet another aspect, the present invention pertains to a device for applying a fastener to suture material attached to tissue within the body of a patient and includes a catheter body having a proximal end and a distal end, an inner body defining a suture recess and an actuation recess, and a movable sleeve defining a deployment lumen. The suture recess on the inner body is in communication with a fastener lumen capable of receiving a fastener therein. The actuation recess is in communication with an actuation lumen formed in the inner body. The deployment lumen formed in the movable sleeve is sized to receive the inner body therein and includes a cutting recess having a cutting member located proximate thereto. [0018] In another aspect, the present invention is directed to a fastener attachable to suture material and comprises a fastener body having at least one attachment lumen formed therein and at least one engagement member attached to the fastener body wherein the engagement member is capable of engaging and retaining the suture material. The engagement member defines an engagement aperture which is in communication with the attachment lumen. The attachment lumen is capable of receiving at least one suture therein. [0019] The present invention also discloses various methods of repairing heart valve tissue within the body of a patient. In one aspect, a method of repairing tissue within the heart of a patient is disclosed which includes advancing a guide catheter through a circulatory pathway to a location in the heart proximate to a heart valve, advancing a therapy catheter through the guide catheter to the heart valve, stabilizing a first leaflet with the therapy catheter, deploying a first suture into the stabilized first leaflet, disengaging the first leaflet from the therapy catheter while leaving the first suture attached thereto, stabilizing a second leaflet with the therapy catheter, deploying a second suture into the second leaflet, disengaging the second leaflet from the therapy catheter while leaving the second suture attached thereto, and joining the first and second leaflets by reducing the distance between the first and second sutures. [0020] An alternate method of repairing tissue within the heart of a patient is disclosed and comprises advancing a guide catheter through a circulatory pathway to a location in the heart proximate to a heart valve, advancing a therapy catheter through the guide catheter to the heart valve, stabilizing a first leaflet with the therapy catheter, deploying a first suture into the stabilized first leaflet, disengaging the first leaflet from the therapy catheter while leaving the first suture attached thereto, stabilizing a second leaflet with said therapy catheter, deploying a second suture into the second leaflet, disengaging the second leaflet from the therapy catheter while leaving the second suture attached thereto, and removing the therapy catheter from the guide catheter. Thereafter, a fastener catheter is positioned over the first and second suture and advanced through the guide catheter to the heart valve. Once positioned, the first and second leaflets are joined by reducing the distance between the first and second sutures and a fastener is deployed from the fastener catheter. [0021] Other objects, features, and advantages of the present invention will become apparent from a consideration of the following detailed description. BRIEF DESCRIPTION OF THE DRAWINGS [0022] The apparatus of the present invention will be explained in more detail by way of the accompanying drawings, wherein: [0023] FIG. 1 shows a perspective view of an embodiment of the guide catheter of the present invention; [0024] FIG. 2 shows a cross-sectional view of an embodiment of the guide catheter of the present invention; [0025] FIG. 3 shows a cross-sectional view of an alternate embodiment of the guide catheter of the present invention; [0026] FIG. 4 shows a cross-sectional view of the embodiment of the guide catheter shown in FIG. 3 ; [0027] FIG. 5 shows a perspective view of an embodiment of the therapy catheter of the present invention; [0028] FIG. 6 shows an embodiment of the therapy device handle of the present invention; [0029] FIG. 7 shows an perspective view of an embodiment of the elongated body of the present invention having a suture attachment tip attached thereto; [0030] FIG. 8A shows a cross-sectional view of an embodiment of the elongated body of the present invention; [0031] FIG. 8B shows a cross sectional view of an alternate embodiment of the elongated body of the present invention; [0032] FIG. 8C shows a side cross-sectional view of the embodiment of the elongated body shown in FIG. 8B ; [0033] FIG. 9 shows a top cross-sectional view of an embodiment of the elongated body of the present invention; [0034] FIG. 10 shows a side cross-sectional view of the embodiment of the elongated body shown in FIG. 9 prior to actuation; [0035] FIG. 11 shows a side cross-sectional view of an embodiment of the elongated body shown in FIG. 10 during actuation; [0036] FIG. 12 shows a side cross-sectional view of an embodiment of the elongated body shown in FIG. 10 following actuation; [0037] FIG. 13 shows another side cross-sectional view of an embodiment of the elongated body shown in FIG. 10 during actuation; [0038] FIG. 14 shows another side cross-sectional view of an embodiment of the elongated body shown in FIG. 10 following actuation; [0039] FIG. 15 shows a top cross-sectional view of an alternate embodiment of the elongated body of the present invention; [0040] FIG. 16 shows a side cross-sectional view of the embodiment of the elongated body shown in FIG. 15 prior to actuation; [0041] FIG. 17 shows a side cross-sectional view of an embodiment of the elongated body shown in FIG. 15 during actuation; [0042] FIG. 18 shows a side cross-sectional view of an embodiment of the elongated body shown in FIG. 15 following actuation; [0043] FIG. 19 shows a perspective view of an embodiment of the fastener catheter of the present invention; [0044] FIG. 20 shows an embodiment of the fastener catheter handle of the present invention; [0045] FIGS. 21 a and 21 b show a perspective view of the components of the fastener tip of the present invention; [0046] FIG. 22 shows a perspective view of the fastener tip of the present invention having a fastener attached thereto; [0047] FIG. 23 shows a side view of an embodiment of the fastener of the present invention; [0048] FIG. 24 shows a side view of the fastener of the present invention attached to suture material; [0049] FIG. 25 shows a perspective view of a guidewire traversing the mitral valve within a heart; [0050] FIG. 26 shows a perspective view of a guide catheter positioned proximate to the mitral valve within a heart; [0051] FIG. 27 shows a perspective view of a therapy catheter advancing through a guide catheter to a position proximate to the mitral valve of a heart; [0052] FIG. 28 shows a perspective view of a therapy catheter stabilizing a first leaflet of the mitral valve of a heart; [0053] FIG. 29 shows a perspective view of the first leaflet of the mitral valve having a suture applied thereto; [0054] FIG. 30 shows a perspective view of a therapy catheter stabilizing a second leaflet of the mitral valve of a heart; [0055] FIG. 31 shows a perspective view of the first and second leaflets of the mitral valve having sutures applied thereto; [0056] FIG. 32 shows a perspective view of a fastener catheter advancing through a guide catheter to a position proximate to the mitral valve of a heart; [0057] FIG. 33 shows a perspective view of the fastener catheter of the present invention applying a fastener to suture material attached to the mitral valve; [0058] FIG. 34 shows a perspective view of the fastener applied to suture material attached to the first and second leaflet of the mitral valve; [0059] FIG. 35 shows a perspective view of another embodiment of the present invention wherein a dilator is used to introduce the guide catheter onto the left atrium; [0060] FIG. 36 shows a perspective view of the dilator of the present embodiment traversing the atrial septum; [0061] FIG. 37 shows a perspective view of the guide catheter of the present embodiment positioned within the left atrium proximate to the mitral valve; [0062] FIG. 38 shows a perspective view of an alternate embodiment of the therapy catheter advanced through the guide catheter to the mitral valve; [0063] FIG. 39 shows a perspective view of the embodiment of the therapy catheter shown in FIG. 38 having an inflatable positioning balloon positioned thereon inflated; [0064] FIG. 40 shows a perspective view of the embodiment of the therapy catheter shown in FIG. 38 engaging a first leaflet; [0065] FIG. 41 shows a perspective view of the first leaflet of the mitral valve having a suture attached thereto; [0066] FIG. 42 shows a perspective view of the embodiment of the therapy catheter shown in FIG. 38 engaging the second leaflet of the mitral valve; [0067] FIG. 43 shows a perspective view of the first and second leaflets of the mitral valve having sutures attached thereto; and [0068] FIG. 44 shows another perspective view of the first and second leaflets of the mitral valve having sutures attached thereto. DETAILED DESCRIPTION OF THE INVENTION [0069] Disclosed herein is a detailed description of various embodiments of the present invention. This description is not to be taken in a limiting sense, but is made merely for the purpose of illustrating the general principles of the invention. The overall organization of the detailed description is for the purpose of convenience only and is not intended to limit the present invention. [0070] The mitral valve repair system of the present invention is designed for use in a surgical treatment of bodily tissue. As those skilled in the art will appreciate, the exemplary mitral valve repair system disclosed herein is designed to minimize trauma to the patient before, during, and subsequent to a minimally invasive surgical procedure while providing improved tissue stabilization and enhanced placement of a fastening device thereon. The mitral valve repair system of the present invention includes a guide catheter capable of being introduced into body of a patient and advanced to an area of interest, a therapy catheter capable of traversing or otherwise engaging the guide catheter and applying a suture to a repair site, and a fastener catheter capable of applying a fastening device to the attached suture. While the guide catheter, therapy catheter, and fastener catheter cooperatively enable a surgeon to deliver a suture to a repair site in vivo, the various components of the present invention may be used individually. For example, the therapy catheter, the fastener catheter, or both may be coupled to a guidewire and advanced to a repair site in vivo without the use of the guide catheter. The mitral valve repair system of the present invention is useful in repairing dysfunctional mitral valve tissue by stabilizing discreet valvular tissue pieces and deploying a fastening device therethrough. However, the mitral valve repair system may be used to repair tissue throughout a patient's body as desired. For example, the present invention may also be used to repair arterial septal defects (ASD), ventricular septal defects (VSD), and defects associated with patent foramen ovale (PFO). [0071] FIGS. 1-4 show various illustrations of the guide catheter of the present invention. As shown in FIG. 1 , the guide catheter 10 comprises a guide body 12 having a proximal end 14 and a distal end 16 . Those skilled in the art will appreciate that the guide catheter 10 of the present invention may be manufactured from a variety of materials, including, without limitation, various plastics, thermoplastics, silicones, elastomers, ceramics, composite materials, or various combinations of the aforementioned materials. In addition, the guide catheter 10 may be manufactured in various lengths and widths as desired by the user. FIGS. 2-4 show various embodiments of the guide catheter 10 . As shown in FIG. 2 , the guide catheter 10 includes an outer wall 18 defining at least one internal lumen 20 . FIGS. 3-4 illustrate alternate embodiments wherein the outer wall 18 defines an internal lumen 20 and includes at least one directing lumen 22 formed therein. The directing lumen 22 is sized to receive a guidewire (not shown) or steering device (not shown) therein. In another embodiment, at least one flexible support structure such as a coiled wire support (not shown) may be embedded within the outer wall 18 of the guide catheter 10 . [0072] FIG. 5 shows a perspective view of an embodiment of the therapy catheter 30 of the present invention. As shown in FIG. 5 , the therapy catheter 30 includes an elongated body 32 having a therapy device handle 34 located at the proximal end and a suture attachment tip 36 located at the distal end Like the guide body 12 of the guide catheter 10 , the elongated body 32 may be manufactured in a variety of shape, sizes, lengths, widths, and biologically-compatible materials as desired. [0073] FIG. 6 shows a more detailed illustration of the therapy device handle 34 of the present invention. As shown, the therapy device handle 34 comprises a handle body 38 having at least a suction connector 40 and a elongated body receiver 42 attached thereto. The suction connector 40 is capable of coupling to a vacuum source (not shown). The elongated body receiver 42 is capable of receiving the elongated body 32 ( FIG. 5 ) thereon. A first actuator 44 is located within a first actuator recess 46 formed on the handle body 38 . Similarly, a second actuator 48 is positioned within a second actuator recess 50 formed in the handle body 38 . As shown in FIG. 6 , a suction actuator 52 , configured to open or close the fluid path between suction connector 40 and elongated body receiver 42 , may be located within a suction actuator recess 54 proximal to the first and second actuators 44 , 48 . [0074] FIGS. 7-10 show various illustrations of the elongated body 32 and the suture attachment tip 36 of the present invention. As shown in FIG. 7 , the elongated body 32 includes a suction recess 56 having a first needle port 58 A and a second needle port 58 B located proximate thereto. The elongated body 32 or the suture attachment tip 36 may include a guidewire port 60 capable of receiving a guidewire 62 . FIG. 8A shows a cross sectional view of the elongated body 32 . As shown, the elongated body 32 comprises an outer wall 64 defining a suction lumen 66 . The suction lumen 66 is in fluid communication with the suction recess 56 ( FIG. 7 ) and the vacuum source (not shown) attached to the suction connector 40 located on the therapy device handle 34 ( FIG. 6 ). A first needle lumen 68 having a first needle 70 located therein may be formed in or otherwise positioned proximate to the outer wall 64 of the elongated body 32 . Similarly, a second needle lumen 72 having a second needle 74 located therein may be formed in or otherwise positioned proximate to the outer wall 64 of the elongated body 32 . The first and second needles 70 , 74 are coupled to or otherwise in communication with the first and second actuators 44 , 48 located on the therapy device handle 34 ( FIG. 6 ). The forward and rearward movement of the first and second actuators 44 , 48 results in the longitudinal movement of the first and second needles 70 , 74 thereby permitting the first and second needles, 70 , 74 to extend from and retract into the first and second needle lumens 68 , 72 . Those skilled in the art will appreciate that the first and second needles 70 , 74 may be capable of individual or simultaneous movement. A first suture lumen 76 having a first suture 78 located therein and a second suture lumen 80 having a second suture 82 located therein may be formed within or located proximate to the outer wall 64 of the elongated body 32 . Of course one of skill in the art will recognize that references herein to “sutures” include not just traditional suture material, but also any material of sufficient length and flexibility to accomplish the purposes of this tissue repair system. In one embodiment, a guidewire lumen 84 sized to receive guidewire 62 therein may be positioned within or proximate to the outer wall 64 of the elongated body 32 and may be in communication with the guidewire port 60 formed on the suture attachment tip 36 . [0075] FIGS. 8B-8C show various illustrations of an alternate embodiment of the present invention, wherein an inflatable positioning balloon 252 is positioned on the outer wall 64 of the elongated body 32 . As shown, the inflatable positioning balloon 252 is in fluid communication with an inflation lumen 84 ′ positioned within the elongated body 32 . The inflation lumen 84 ′ may be in fluid communication with an inflation source in ways known to those skilled in the art and may be attached to or otherwise in communication with the therapy device handle 34 ( FIG. 5 ), thereby permitting the position of the therapy catheter 30 to be manipulated without using a guidewire. Moreover, the positioning balloon 252 can be used to hold the therapy device steady once in position. [0076] FIGS. 9-10 show various illustrations of the present invention prior to use. As shown, a first needle receiving port 86 A may be positioned within or proximate to the suction lumen 56 co-aligned with and opposing the first needle port 58 A. Similarly, a second needle receiving port 86 B may be positioned within or proximate to the suction lumen 56 co-aligned with and opposing the second needle port 58 B. The first needle receiving port 86 A is in communication with the first suture lumen 76 and contains at least a first needle catch 88 A attached to the first suture 78 therein. Likewise, the second needle receiving port 86 B is positioned proximate to the suction recess 56 opposing the second needle port 58 B. The second needle receiving port 86 B is in communication with the second suture lumen 80 and contains a second needle catch 88 B attached to the second suture 82 therein. [0077] FIGS. 11-12 show an embodiment of the therapy catheter of the present invention during various stages of use. As shown in FIG. 11 , forward movement of the first actuator 44 within the first actuator recess 46 ( FIG. 6 ) results in the first needle 70 advancing through the first needle port 58 A and traversing the suction recess 56 . Continued actuation of the first actuator 44 results in the first needle 70 advancing through the first needle receiving port 86 A and engaging the first needle catch 88 A positioned within the first suture lumen 76 . The first needle catch 88 A engages and is retained on the first needle 70 . The user may then retract the first needle 70 , thereby pulling the first suture across the suture recess 56 . To retract the first needle 70 , the user rearwardly moves the first actuator 44 . As shown in FIG. 12 , the first needle 70 having the first needle catch 88 A attached thereto is retracted through the first needle receiving port 86 A, traverses the suction recess 56 , and enters the first needle lumen 68 through the first needle port 58 A. FIG. 12 shows the first suture 78 traversing the suction recess 56 . [0078] Similarly, as shown in FIG. 13 , forward movement of the second actuator 48 ( FIG. 6 ) results in the second needle 74 advancing through exiting the second needle port 58 B and traversing the suction recess 56 . Like the actuation process described above, the continued actuation of the second actuator 48 results in the second needle 74 advancing through the second needle receiving port 86 B and engaging the second needle catch 88 B positioned within the second suture lumen 80 . The second needle catch 88 B is then engaged and retained on the second needle 74 . Thereafter, the user may retract the second needle 74 thereby pulling the second suture across suture recess second needle port 58 B. To retract the second needle 74 , the user rearwardly moves the second actuator 48 . As shown in FIG. 14 , the second needle 74 having the second needle catch 88 B attached thereto is retracted through the second needle receiving port 86 B, traverses the suction recess 56 , and enters the second needle lumen 72 through the second needle port 58 B. The second suture 82 , which is attached to the second needle catch 88 B, thus traverses the suction recess 56 . [0079] FIG. 15 illustrates an alternate embodiment of the present invention. As shown, the elongated body 32 includes a suction recess 90 formed thereon which is in fluid communication with a suction lumen 92 formed therein which in turn is in communication with a vacuum source (not shown) attached to the suction connector 40 ( FIG. 6 ). First and second needle ports 94 A, 94 B, respectively, are positioned within or proximate to the suction recess 90 . Similarly, first and second needle receiving ports 96 A, 96 B, respectively, are positioned within or proximate to the suction recess 90 and are co-aligned with and opposed to the first and second needle ports 94 A, 94 B. The first needle port 94 A communicates with a first needle lumen 98 . A first deployment rod 100 having a first detachable needle 102 attached thereto is located within the first needle lumen 98 . The first detachable needle 102 is coupled to a first suture 104 located within the first needle lumen 98 . Similarly, the second needle port 94 B communicates with a second needle lumen 106 . A second deployment rod 108 having a second detachable needle 110 attached thereto is located within the second needle lumen 106 . The second detachable needle 110 is coupled to a second suture 112 located within the second needle lumen 106 . The first needle receiving port 96 A leads to a first needle trap lumen 114 A formed in or positioned proximate to suction recess 90 . A first needle trap 116 A capable of receiving and retaining the first detachable needle 102 therein is positioned within the first needle trap lumen 114 A. Similarly, the second needle receiving port 96 B leads to a second needle trap lumen 114 B formed in or positioned proximate to the suction recess 90 Like the first needle trap 116 A, a second needle trap 116 B capable of receiving and retaining the second detachable needle 110 therein is positioned within the second needle trap lumen 114 B. [0080] FIGS. 16-18 show the embodiment of FIG. 15 during use. Forward movement of the first actuator 44 results in first needle rod 100 extending from first needle lumen 98 . FIG. 17 shows the first needle rod 100 with a first detachable needle 102 attached thereto extended through the first needle port 94 A traversing the suction recess 90 , and entering into the first needle trap lumen 114 A through the first needle receiving port 96 A. The first detachable needle then engages the first needle trap 116 A. Thereafter, the first needle rod 100 is retracted into the first needle lumen 98 , thereby leaving first detachable needle 102 in first needle trap 116 A. To retract the first needle rod 100 , the user moves the first actuator 44 a rearward direction which causes the first needle rod 100 to retract into the first needle lumen 98 . FIG. 18 shows the first needle rod 100 retracted into the first needle lumen 98 . As a result, the first suture 104 which is attached to the first detachable needle 102 traverses the suction recess 90 . Those skilled in the art will appreciate that a second needle (not shown) may be deployed in a similar manner. [0081] FIGS. 19-21 show various illustrations of the fastener catheter of the present invention. As shown in FIG. 19 , the fastener catheter 130 comprises a fastener catheter body 132 having a fastener catheter handle 134 attached at the proximal end and a fastening tip 136 at the distal end. The fastener catheter 130 may be manufactured in a variety of shapes, sizes, lengths, widths, and biologically-compatible materials as desired. [0082] FIG. 20 shows a more detailed illustration of a preferred fastener catheter handle 134 of the present invention. As shown, the fastener catheter handle 134 comprises a fastener handle body 138 having an auxiliary connector 140 and a fastener body receiver 142 attached thereto. The auxiliary connector 140 may be capable of coupling to a variety of devices including, for example, a vacuum source or a visualization device. The fastener body connector 142 is capable of receiving and coupling to the fastener catheter body 132 ( FIG. 19 ). A fastener actuator 144 may be positioned within a fastener actuator recess 146 formed on the fastener handle body 138 . The fastener actuator 144 positioned within the fastener actuator recess 146 may be capable of being positioned in three distinct locations. For example, in a non-actuated condition, the fastener actuator 144 may be located in a first position 148 . Thereafter, the user may partially actuate the fastener catheter 130 by positioning the fastener actuator 144 in a second position 150 , thereby deploying a fastening device (not shown) from the fastener catheter 130 ( FIG. 19 ). The user may then fully actuate the fastener catheter 130 by moving the fastener actuator 144 to a third position 152 within the fastener actuator recess 146 , thereby actuating a cutting member (discussed below) located on or proximate to the fastening tip 136 . [0083] FIGS. 21 a and 21 b illustrate, in exploded fashion, pieces of fastening tip 136 . An inner body 154 includes a suture recess 160 formed in the side thereof, which in turn is in communication with an internal fastener lumen 158 . Inner body 154 also includes a pin 162 extending radially outward therefrom. Sleeve 156 comprises an axial deployment lumen 166 of sufficient diameter to receive inner body 154 therein. Sleeve 156 also comprises a cutting recess 168 formed in an axial side thereof and a cutting member 170 on a proximal edge of cutting recess 168 . Slot 172 extends parallel to the axis of the deployment lumen 166 and may extend radially through to fastener lumen. Pin recess 172 receives pin 162 in sliding relation. [0084] FIGS. 23-24 illustrate fastener 180 of the present invention. Fastener 180 may be manufactured from a variety of materials including, for example, Nickel-Titanium alloys, shape-memory alloys, stainless steel, titanium, various plastics, and other biologically-compatible materials. Fastener 180 has an internal lumen 188 extending axially therethrough and one or more engagement member(s) 184 formed on an end thereof. Between the engagement members is defined engagement aperture 186 which is in communication with attachment lumen 188 . Attachment lumen 188 and engagement aperture 186 are sized to receive a first suture lead 176 A and a second suture lead 176 B therein. Prior to deployment, engagement member(s) 184 are deflected radially away from the axis of the fastener 180 such that engagement aperture 186 has a relative large first diameter sufficient to permit suture leads 176 A and 176 B to slide therethrough. Upon deployment, i.e. after the suture leads 176 A and 176 B have been retracted, engagement members 184 are deflected or permitted to spring back toward the axis of the device such that the engagement aperture 186 assumes a second smaller diameter compressing and securing suture leads 176 A and 176 B in place. Preferably the engagement member(s) 184 tend to spring toward a natural position at the axis of fastener 180 . FIG. 24 shows the fastener 180 in the deployed configuration in which a suture loop 178 has passed through two discreet tissue portions 200 A, 200 B and suture leads 176 A, 176 B are secured in fastener 180 . Each engagement member(s) 184 may further include a pointed tip 190 which, when the engagement member(s) are in the deployed position, engages and further restricts movement of the suture leads 176 A, 176 B. [0085] An operational fastening tip 136 with fastener 180 attached thereto and ready for deployment can be seen in FIG. 22 . Inner body 154 has been placed inside sleeve 156 such that suture recess 160 is in alignment with cutting recess 168 . Pin 172 is in slidable communication with slot 162 thereby permitting relative linear motion, but preventing relative rotational motion, between inner body 154 and sleeve 156 . Fastener 180 has been placed on the end of the fastening tip 136 by deflecting the engagement members 184 radially outward until they can be placed around the outer circumference of the inner body 154 . Accordingly, the fastener is secured to the end of inner body 154 by means of the frictional engagement between the engagement members 184 and the outer surface of inner body 154 . Suture loop 178 extends from the fastener 180 . Suture leads 176 A and 176 B extend through the lumen 188 , through engagement aperture 186 , exit the side of inner body 154 through suture recess 160 , and exit the side of sleeve 156 through cutting recess 168 . [0086] Deployment of the fastener is a two step process. Once suture 178 has been secured through one or more tissue segments, the fastener tip 136 is coaxed toward the tissue and the suture leads 176 A and 176 B are pulled away from the tissue until the suture loop is sufficiently cinched around the target tissue. Sleeve 156 is then held in place adjacent the tissue while the inner body 154 is pulled axially away. This causes sleeve 156 to push (i.e. slide) fastener 180 off the outer surface of the inner body 154 . When fastener 180 has been completely removed from inner body 154 engagement members 184 spring axially inward thereby reducing the diameter of engagement aperture 186 and securing suture leads 176 A and 176 B. The second deployment step, cutting suture leads 176 A and 176 B, is accomplished when the inner body 154 is pulled sufficiently through sleeve 156 that the suture leads are pinched between the trailing edge of suture recess 160 and cutting member 170 and ultimately cut by cutting member 170 . [0087] Remote deployment of fastener 180 is accomplished by attaching inner body 154 to fastener actuator 144 , and attaching sleeve 156 to the fastener catheter handle 134 . Thus, axial movement of the fastener actuator 144 relative to the handle 134 causes similar relative movement between inner body 154 and sleeve 156 . For example, in the non-actuated position 148 (see FIG. 20 ) the distal end of inner body 154 will extend from sleeve 156 a sufficient distance to hold fastener 180 thereon. In the second position 150 the inner body 154 will have been withdrawn into sleeve 156 a sufficient distance to deploy the fastener 180 , and in the third position 152 the inner body 154 will have been withdrawn a sufficient distance to cut the suture leads 176 A and 176 B. [0088] The present invention also discloses various methods of using the disclose mitral valve repair system to repair discreet tissue portions in vivo. The following paragraphs describe methods of repairing a dysfunctional mitral valve, though those skilled in the art will appreciate that the present invention and procedure may be adapted for use on other valves or in other procedures requiring the attachment of two or more pieces of tissue. [0089] To repair a dysfunctional or otherwise incompetent heart valve, a guidewire capable of traversing the circulatory system and entering the heart of the patient is introduced into the patient through an endoluminal entry point. For example, the endoluminal entry point may be formed in a femoral vein or right jugular vein. Thereafter, the guidewire is advanced through the circulatory system, eventually arriving at the heart. The guidewire is directed into the right atrium, traverses the right atrium and is made to puncture with the aid of a tran-septal needle or pre-existing hole, the atrial septum, thereby entering the left atrium. As shown in FIG. 25 , the guidewire 220 may then be advanced through the mitral valve 222 and into the left ventricle 226 . The guidewire 220 traverses the aortic valve 228 into the aorta 230 and is made to emerge at the left femoral artery through an endoluminal exit point. Once the guidewire 220 is positioned, the endoluminal entry or exit port is dilated to permit entry of a catheter therethrough. A protective sheath may be advanced in the venous area to protect the vascular structure. [0090] As shown in FIG. 26 , the guide catheter 10 of the present invention may be attached to the guidewire 220 and advanced through the dilated guidewire entry port to a point proximate to the mitral valve 222 . Those skilled in the art will appreciate that the mitral valve repair system of the present invention may approach the mitral valve from an antegrade position or from a retrograde position as desired by the user. Once the guide catheter is suitably positioned in the heart, the therapy catheter 30 may be advanced through the guide catheter 10 to a position proximate to the mitral valve 222 . FIG. 27 shows the therapy catheter 30 emerging from the guide catheter 10 proximate to the mitral valve 222 . Thereafter, the user may actuate the suction actuator 52 located on the handle body 38 of the therapy device handle 34 ( FIG. 6 ). As a result, a suction force is applied from the suction recess 56 formed on the suture attachment tip 36 of the therapy catheter 30 ( FIG. 7 ) to the tissue located proximate thereto. As shown in FIG. 28 , a first valve leaflet 240 A is engaged and retained by the suction force applied through the suction recess 56 . With the first valve leaflet 240 A stabilized, the user may apply a suture 242 A thereto as described above. To apply the first suture to the first valve leaflet 240 A, the user actuates the first actuator 44 located on the therapy device handle 34 , which results in the first needle 70 advancing through the first valve leaflet 240 A and engaging and retaining the first needle catch 88 A, thereby applying a first suture 242 A to the tissue ( FIGS. 6-7 ). Thereafter, the user may terminate application of suction force to the first valve leaflet 240 A thereby releasing the sutured tissue. FIG. 29 shows the first valve leaflet 240 A having a first suture 242 A applied thereto. As shown in FIG. 30 , the therapy catheter 30 may then be rotated and positioned to engage a second valve leaflet 240 B. Once again, the user may actuate the suction actuator 52 to apply suction force to the second valve leaflet 240 B through the suction recess 56 . With the second valve leaflet 240 B stabilized as shown in FIG. 30 , the user may apply a suture 242 B thereto by actuating the second actuator 48 located on the therapy device handle 34 , which results in the second needle 74 advancing through the second valve leaflet 240 B and engaging and retaining the second needle catch 88 B, thereby applying a second suture 242 B to the tissue. As shown in FIG. 31 , the user may terminate the application of suction to the stabilized tissue and remove the therapy catheter from the patient, thereby leaving the first and second sutures 242 A, 242 B attached to the first and second valve leaflets 240 A, 240 B. Note that first and second sutures 242 A and 242 B are actually portions of the same suture such that when the therapy catheter is removed there is a single suture loop through the valve leaflets 240 A and 240 B. [0091] As shown in FIGS. 32-33 , the fastener catheter 130 may be attached to the guidewire 220 and will be attached to first and second sutures 242 A, 242 B. Thereafter, the fastener catheter 130 may be inserted into the guide catheter 10 and advanced to a position proximate to the mitral valve 222 . The user then draws the first and second sutures 242 A, 242 B taut while advancing the fastener catheter 130 to the mitral valve 22 , thereby decreasing the distance between the first and second valve leaflets 240 A, 240 B. The user then actuates the fastener actuator 144 which causes the sleeve 156 to engage and apply the fastener 180 to the first and second sutures 242 A, 242 B adjacent the leaflets, as described above. Continued actuation of the fastener actuator 144 causes the cutting member 170 to engage and cut the first and second sutures 242 A, 242 B. As shown in FIG. 34 , after the fastener catheter 130 , the guide catheter 10 , and the guidewire 220 are removed from the patient, the fastener 180 remains applied to the mitral valve 222 . [0092] FIGS. 35-44 describe an alternate method of repairing tissue, specifically valve leaflets in this embodiment, in vivo. As shown in FIG. 35-37 , a guide catheter 10 is advanced through the circulatory system to the right atrium of the heart. Once positioned, a dilator 250 is advanced through the guide catheter 10 and is made to puncture the atrial septum, thereby entering the left atrium. Thereafter, the guide catheter 10 is advanced into the left atrium through the punctured atrial septum and positioned proximate to the mitral valve 222 . As shown in FIG. 38 , the therapy catheter 30 may be inserted into the guide catheter 10 and advanced to a position proximate to the mitral valve 222 . As shown in FIG. 39 , an inflatable positioning balloon 252 (discussed above) located on the therapy catheter 30 is inflated to orient and steady the catheter. The suction actuator 52 on the therapy device handle 34 is then actuated to apply a suction force to the suction recess 56 (see. FIG. 6 ). The inflated balloon 252 engages the second valve leaflet 240 B which forces the suction recess 56 towards the first valve leaflet 240 A, thereby resulting in the stabilization of the first valve leaflet 240 A as shown in FIG. 40 . As shown in FIG. 41 , the user may then apply the first suture 242 A to the first valve leaflet 240 A as described above. Once the suture is applied, the user may deflate the inflatable positioning balloon 252 and rotates the therapy catheter 30 approximately 180 ° . Thereafter, the user inflates the positioning balloon 252 and actuates suction actuator 52 to apply a suction force to the suction recess 56 . As shown in FIG. 42 , the inflatable positioning balloon 252 is again inflated and made to engage the first valve leaflet 240 thereby forcing the suction recess 56 to engage the second valve leaflet 240 B and permitting the stabilization of the second valve leaflet 240 B as shown in FIG. 42 . Thereafter, the user applies the second suture 242 B to the second valve leaflet 240 B as described above. FIGS. 43-44 show the first and second valve leaflets 240 A, 240 B having a first and second suture 242 A, 242 B applied thereto. Thereafter, the therapy catheter 30 is removed from the patient's body and the fastener catheter 130 is used to apply a fastener to the first and second sutures 242 A, 242 B as described above. [0093] In closing, it is understood that the embodiments of the invention disclosed herein are illustrative of the principals of the invention. Other modifications may be employed which remain within the scope of the present invention. Accordingly, the present invention is not limited to the embodiments shown and described in this disclosure.

The present invention is directed to a method of securing suture within a patient. The method comprises advancing suture material into the patient, passing a portion of the suture material through tissue of the patient, and advancing a fastener catheter and a suture fastener along the suture material. The fastener catheter includes a catheter main body and a handle. The suture fastener has a generally cylindrical body formed of a shape memory material and a suture fastener inner lumen. The suture fastener has at least one engagement tab biased to extend at least partially into the suture fastener inner lumen to engage against and secure one or more lines of suture passing through the suture fastener.

BACKGROUND OF THE INVENTION 1. Technical Field The present invention relates to broadband signal distribution equipment for distributing communication signals, and, in particular, to a method and apparatus of broadband signal distribution whereby interference with the broadband signal distribution can be prevented when access to an enclosure housing the distribution equipment is authorized. 2. Description of the Prior Art The theft of cable television signals by unauthorized persons has plagued the cable television industry since its inception. New systems and methods are continually being developed to prevent signal theft. One category of systems for preventing signal theft utilizes a signal distribution apparatus for distributing CATV signals locally to various subscribers. Such a distribution apparatus is disclosed in U.S. Pat. No. 4,963,966 to Harney et al. Piracy is prevented by encasing the distribution apparatus in an enclosure and locating the enclosure off of the subscribers' premises. As a result, pirates have reduced opportunity to examine and effect changes to system circuitry of the signal distribution apparatus. The enclosures are capable of both aerial mounting on, for example, a telephone pole, or pedestal mounting for use at installations of underground cable or for placement in a locked room, for example. In addition, the enclosures contain tamper prevention circuitry for disabling the signal distribution apparatus when the enclosure is opened. When service of the signal distribution apparatus is necessary, the tamper prevention circuitry can be inhibited through a signal transmitted from the headend from which the broadband signal originates. A serviceman can then access the enclosure without disabling the signal distribution apparatus. The signal distribution apparatus can perform such functions as premium channel interdiction in which an interfering signal is introduced into the television signal at a subscriber's location. This technique ensures that events carried on premium channels are available only to those subscribers authorized to receive the events. Such an interdiction system is described in U.S. Pat. No. 4,912,760 to West, Jr. et al. The signal distribution apparatus contains at least one microprocessor controlled oscillator and switch control electronics to secure several premium television channels. Control is accomplished by injecting an interfering or jamming signal into unauthorized channels from the pole-mounted unit. To improve efficiency and to save costs, one oscillator may be used to jam several premium television channels. This technique reduces the amount of required hardware and maximizes system flexibility. The oscillator output jamming signal frequency is periodically moved from channel to channel. Consequently, the oscillator is frequency agile and hops from jamming one premium channel frequency to the next. Costs are reduced since a single interdiction unit may serve a plurality of subscribers. The cable television industry has additionally been pressured by subscribers to service defective system equipment promptly in order to minimize periods of signal loss. Accordingly, service of system equipment, including distribution apparatus, must be conducted around the clock. In order to reduce the costs of providing such comprehensive service, the number of employees and the amount of time needed to access and service a particular distribution apparatus must be reduced. However, if access is made easier, the distribution apparatus becomes susceptible to piracy. By the same token, systems and methods for preventing signal theft oftentimes make the service of distribution apparatus slow, inconvenient, and expensive. Accordingly, there is a desire in the industry to provide a signal distribution apparatus that prevents piracy but is still easily and efficiently serviced. SUMMARY OF THE INVENTION It is an object o the present invention to provide a broadband signal distribution apparatus which prevents piracy and can be efficiently serviced. It is a further object of the present invention to provide a broadband signal apparatus having a tamper prevention system which may be serviced without the need for communication with and assistance from the headend. It is a further object of the present invention to provide a tamper override system in a broadband signal distribution apparatus which is secure from pirates, is efficiently serviced, and does not require significant additional expense for the implementation. It is further object of the invention to provide a method for screening access to a broadband signal distribution system which both prevents signal theft and facilitates access by authorized personnel. It is a further object of the present invention to provide a tamper override module for locally overriding a tamper prevention system contained in a signal distribution system without the need for communication with and assistance from the headend. In accordance with the present invention, a broadband signal distribution apparatus includes a signal receiver for receiving a broadband signal from a headend and a signal distributor for distributing the broadband signal to at least one subscriber location. The signal receiver and the signal distributor are housed in an enclosure. The signal distribution apparatus additionally includes a tamper for detecting when the enclosure is opened, a device for applying an entry code, and a storage device for storing an authorized access code. A comparator compares the applied entry code with the authorized access code. In response to the tamper detector and the comparator, interfering circuitry interferes with the operation of the signal distributor if the enclosure is opened and the entry code does not match the authorized access code. The present invention achieves the objective of preventing piracy by providing a secure enclosure, a tamper detector, and interfering circuitry for interfering with the operation of the signal distribution apparatus when the enclosure is opened by an unauthorized person. As a result, operator revenue is not depleted by service pirates. In addition, the present apparatus achieves the objective of facilitating service of the signal distribution apparatus by enabling an authorized serviceman to apply an entry code which prevents operation of the interfering circuitry. Accordingly, additional personnel are not required to assist a serviceman in accessing the signal distribution apparatus. Furthermore, reliance upon communication between a local serviceman and personnel at the headend is not required. Moreover, the invention may be realized without significant additional expenditures for new equipment and installation. Accordingly, the present invention prevents piracy while permitting service to be accomplished rapidly and at a reduced cost. In further accordance with the present invention, a method is provided for screening access to a broadband signal distribution apparatus including a broadband signal receiver for receiving a broadband signal from a headend, a broadband signal distributor coupled to said broadband signal receiver for distributing the broadband signal to at least one subscriber location, and an enclosure housing the broadband signal receiver and broadband signal distributor. The method includes the steps of storing an authorized access code, detecting access to the enclosure, and determining whether an entry code is applied, the entry code is compared to the authorized access code. The method further includes the step of interfering with the broadband signal distributor if an entry code is not applied or if an applied entry code does not match the authorized access code. The present invention achieves the objective of preventing signal theft by interfering with the operation of the broadband signal distributor. However, the interfering step can be avoided if an entry code which matches a stored authorized access code is applied. Therefore, access to the broadband signal distribution apparatus by authorized personnel may be accomplished. Furthermore, the present invention provides a tamper override module including a device for setting an entry code, an interface adapted to be connected to a broadband signal distribution apparatus, and an output device for outputting the entry code through the interface to the broadband signal distribution apparatus. The present invention achieves the objective or overriding a tamper prevention system without the need for communication with and assistance from the headend by providing a tamper override module which may be connected to a broadband signal distribution apparatus. The tamper override module outputs an entry code to the broadband signal distribution apparatus. Circuitry in the broadband signal distribution apparatus prevents operation of the tamper prevention system in response to the entry code. Accordingly, access to the broadband signal distribution apparatus is permitted without requiring further assistance from or communication with the headend. Consequently, subscribers are provided the service they demand, servicemen can perform their job without unnecessary delay, and cable operators are protected from pirates without significant additional expense. The above and further objects and advantages of the invention will become apparent with reference to the detailed disclosure of the invention below and the accompanying illustrative figures. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 illustrates a cable television distribution system including a plurality of signal distribution apparatuses. FIG. 2 is a cutaway view of a broadband signal distribution apparatus having an interface with a tamper override module coupled thereto. FIG. 3 is a view from above of a broadband signal distribution apparatus having an interface with a tamper override module coupled thereto. FIG. 4 illustrates a perspective view of a tamper override module. FIG. 5 is block schematic diagram of an addressable common control circuit for a plurality of provided subscriber modules of a signal distribution apparatus including a broadband signal tap, a diplexer connected to the tap, a microprocessor coupled to a tamper override module, tamper detector, a data receiver and decoder, and an automatic gain circuit. FIG. 6 is a block schematic diagram of one subscriber module of a signal distribution apparatus comprising a microprocessor for selectively controlling the jamming of unauthorized services to a subscriber, associated jamming equipment and a diplexer. FIG. 7 is a schematic of a tamper override module. FIG. 8 is a block diagram illustrating the operation of a distribution apparatus and a tamper override module when coupled to an interface of the distribution apparatus. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS A method and apparatus for broadband signal distribution according to the present invention will be discussed in the context of the off-premises cable television distribution apparatus disclosed in U.S. Pat. No. 4,912,760, the disclosure of which is herein incorporated by reference in respect to those features not described in the present specification. However, the present invention is not limited to interdiction cable television systems but is also applicable to other off-premises systems requiring periodic service and operational security. The technology may also be applied to on-premises systems and technologies, for example on-premises addressable converters and addressable descrambler converter technologies. A detailed discussion of the interdiction system in which the present invention may be implemented is also provided in U.S. Pat. No. 4,963,966, also incorporated herein by reference. Furthermore, U.S. patent application Ser. Nos. 446,695; 498,084, 503,423; and 625,901, entitled "CATV Pay Per View System Method and Apparatus," also incorporated herein by reference, disclose details of off-premises subscription television apparatus related by subject matter to the present invention. Topics relating to interdiction systems such as jamming signal frequency calibration, gain adjustment and jamming control oscillator will not be addressed herein in detail. FIGS. 1-8 comprise drawings of the present invention in which similar reference numerals are used throughout to designate similar elements. FIG. 1 is a general block diagram of a subscription television system and, in particular, a cable television system. Subscription television system is intended herein to include any television service system such as over-the-air systems, satellite signal access systems and such television delivery systems. Cable television system as used herein refers to all systems involving the transmission of broadband signals over a transmission medium (fiber optic cable or coaxial cable) to remote locations. For example, a cable television system may comprise a community antenna television distribution system, a satellite signal distribution system, a broadcast television system, a private cable distribution network, either industrial or educational, or other forms of such systems. Each remote location of a television receiver may comprise the location of a particular subscriber to a subscription television service, plural subscribers, single subscribers having plural television receivers or private locations in a private cable distribution network. Consequently, the term subscriber as used herein refers to either a private subscriber or a commercial user of the cable television system. Headend 10 is a connecting point to a serving cable or trunk 28 for distributing, for example, television or other communication channels from service source 12 over feeder lines to drops 66, 88 and finally to subscriber locations. For reference purposes, an Electronic Industries Association (E.I.A.) standard cable television frequency allocation scheme is employed and referred to herein. Typically, in such systems, television channels of source(s) 12 are modulated and frequency division multiplexed together to comprise a broadband signal which may have a bandwidth in excess of one gigahertz. However, one may apply the principles of the present invention to other known standard or non-standard frequency allocations. Further, a National Television Subcommittee (N.T.S.C.) standard television signal at baseband is generally considered in the following description. However, the principles of the present invention apply equally to other standard and non-standard baseband standard definition and proposed high definition television signal formats. Also, the principles of the present invention are not limited to television services furnished from a headend but may include utility meter reading, burglar alarm reporting, digital or other stereophonic audio delivery systems, video or telephonic services and the like. Headend 10 typically comprises a source of television programming 12. The television program source 12 may be a satellite television receiver output, a program produced by a television studio, program material received over a microwave or broadcast television link, a cable television link output, or any other source of television or other subscription programming consistent with the present invention. The program source material need not be limited to conventional television but may comprise teletext, videotext, program audio, utility data, or other forms of communication to be delivered to a remote location over the serving cable or trunk line 28 and subsequently over feeder lines and, then, drop lines 66, 88. Communications used to authorize or terminate services or different levels of service or to perform or monitor system functions are initiated via data entry to a computer system including a system manager computer 16 and a billing computer 18. Conventionally, trunk line 28, feeder lines, and drop lines 66, 88 are constructed of coaxial cable. For higher performance, any one of these lines could be a fiber optic cable. Due to the cost of the installation and the need for a high quality initial transmission from headend 10, trunk line 28 is typically the only line constructed of fiber optic cable. Program material provided by the service source 12 may be included within a regular service offering, or premium (requiring extra payment) or otherwise restricted or desirably secured from receipt at unauthorized receiver locations. It may be provided over any channel of the 50-550 MHz (or larger band) cable television spectrum. "Premium channel" or "premium programming" as used herein refers to a channel or program which is secured from unauthorized receipt either because of its premium or restricted status or because a regular service subscriber must incur an additional fee for reception. Normally, all premium programming in cable television systems is scrambled. However, in accordance with interdiction system technology, premium programming is transmitted in the clear, and interdiction (service denial) is applied at off-premises signal distribution apparatus 20 to jam reception of unauthorized premium programming. Off-premises signal distribution apparatus 20, 22, 24, and 26 may also form part of a reverse data transmission path as will be discussed in greater detail below. Off-premises signal distribution apparatus 20, 22, 24, and 26 respectively comprise housings or enclosures including casings 56, 58, 60, and 62 and covers 48, 50, 52, and 54. It is likely that cable systems will gradually make the transition to an interdiction system, for example, as new subscribers are added. During a transition period, headend 10 may provide scrambled television programming as well as premium programming in the clear and scrambler may be provided as long as converters/decoders remain in the system for unscrambling scramble program transmission. For example, off-premises distribution apparatus 24 may be coupled to subscribers still having on-premises converters/decoders and off-premises unit 22 may be utilized to couple new subscribers to the system. In certain instances, converter/decoders at subscriber locations may later be entirely replaced by interdiction apparatus of the present invention. Descrambling or decoding equipment may also be provided at an off-premises housing. Headend 10 includes an addressable data transmitter 14 for transmitting global commands and data downstream to all subscribers, group addressed communications to a particular group of subscribers, or specifically addressed communications for reception by a particular subscriber. Such forward data transmission may be conducted over a separate data carrier from the cable television spectrum, for example, at 108.2 megahertz. Forward data transmission may also be over an unused default channel from the television spectrum. Global commands generally take the form of operation code and data while addressed communications further include the unique address of a particular subscriber. A high speed data transmission format may be provided for communication between headend 10 and apparatus 56 via transmitter 14. One such format may be a biphase data encoding at a data rate of as least 12 to 14 kilobits per second and preferably 19.2 kilobits per second, frequency modulated on the 108.2 MHz data carrier (in the FM broadcast band). Frequency shift keying, period shift keying, or other related data transmission schemes may be used in the alternative. Principles surrounding appropriate data rate and format selection include meeting objective of handling plural serial communications arrangements and maximizing data throughput. For example, the preferred addressable data transmitter queues communications for transmission and is capable of transmitting hundreds of communications per second. In an alternative embodiment, forward data communications may take the form of in-band signals sent with a television channel superimposed, for example, upon an audio carrier during a special time period, for example, a period corresponding to the vertical blanking interval of the associated video signal, Such data communications further complicate data reception at interdiction apparatus and are desirably eliminated. However, in-band signaling may be required for the operation of certain converter/decoders known in the art. Thus, communications, in particular, commands to authorize service to a particular subscriber, may be transmitted in-band or on a separate data carrier and typically involve transmitting a unique address of a particular subscriber unit or group of subscribers, a command or operation code and data. Decoders in the system receive the command, decode it, determine if the command is to be acted on, and if so perform the desired action such as provide a subscriber with pay-per-view credits or generally authorize services. Responsive to the control of the system manager computer 16, channel program or authorization data is transmitted via and addressable data transmitter 14 over a trunk line 28 to feeder lines with interspersed signal amplifiers 34 and power supply equipment 41, 42, 44, 46 provided as required. The serving signal is dropped via drops 66, 88 to a subscriber location at a pole 36 or from a pedestal at underground cable locations or in equipment closets. Signal distribution apparatus 24 may be connected via connector 86 and drop 88 to a conventional converter/decoder which serves several functions. Responsive to an addressed communication from headend addressable data transmitter 14, channel or program authorization data is updated in an authorization, memory if the address associated with the addressed communication matches a unique address of the subscriber decoder. For example, the subscriber address may comprise a plurality of bits over and above the actual number of subscribers in a system, the additional bits insuring the security of the address. The premium channel or program is then stored in the authorization memory of the converter/decoder. Television programming is normally converted to an otherwise unused channel such as channel 3 or 4 of the television spectrum by a converter portion of converter/decoder. Its premium status is checked against the data stored in authorization memory. If the programming is authorized, the decoder portion of the converter/decoder is enable to decode authorized scrambled premium programming. The provided television receiver TV may be a conventional television receiver or may be a so-called cable ready television receiver. Because of the advent of cable ready television receivers, there is no longer a requirement at a subscriber premises for the converter portion of a converter/decoder because a converter is built into such television receivers. The television receivers may also comprise video cassette recorders (VCRs) or other recording devices which are likewise cable ready and adapted to receiver a signal comprising a regular subscription service offering. A television receiver display may be adaptedly connected by the subscriber to receive over-the-air broadcasts, satellite repeated signals, and other alternative sources of signals such as taped programs via VCRs. Subscriber premises shown at the end of the drops 66, 68 may comprise single family homes, multiple family dwellings such as apartment complexes, hotel, hospitals and the like, or commercial establishments such as restaurants, bars, theaters, factories and the like. The subscriber premises should not be narrowly construed to comprise only single family dwelling units shown in the accompanying drawings. In accordance with a cable television system provided with interdiction or other signal distribution apparatus, units 20, 22, 24, and 26 are mounted on a strand 38 supporting the cable to a pole 36, or provided via a pedestal, as is shown more particularly in U.S. Pat. No. 4,963,966. The units may also be mounted indoors in an equipment closet of a multiple dwelling unit or to the side of a subscriber's premises. Inside the units is common control circuitry for tapping into the broadband television and data transmission spectrum. Referring to the pole 36, there is shown a stand-mounted apparatus 56 serving four drops 66 to subscribers via connector 64. In practice, four or more subscribers and up to four or more drops 66 may be served by signal distribution apparatus 20. In addition to the common control circuitry, four or more plug-in subscriber modules may be provided for an off-premises housing. Also, according to the present invention, additional services requiring two way data transmission such as subscriber polling, home shopping, burglar alarm, energy management and pay-per-view services may be provided via four or more special service modules comprising reverse path signal combining circuitry of apparatus 56. Desirably, all cable television equipment may be removed from the subscriber premises. However, for the provision of certain additional services, some on-premises equipment is unavoidable. For example, a subscriber transaction terminal apparatus may be provided in a subscriber's premises. The subscriber transaction terminal may simply comprise a subscriber-controlled data transmitter for transmitting data on the subscriber drop 66 in only one direction, namely, to signal distribution apparatus 20. For purposes of this description, the subscriber premises will be assumed to include at least one cable ready conventional television receiver, TV or VCR. Consequently, subscriber equipment need not comprise a tunable converter for converting a received cable television channel to an unused channel such as channel 3 or 4. The subscriber transaction terminal device comprises data entry or sensing means, data confirmation means, i.e., a display or alarm, if required, and a data transmitter coupled between the drop cable and the cable ready television receiver. Power for signal distribution apparatus 20 may be provided over the cable from the headend direction via power supplies 41,42 or be provided via the subscriber drip 66 or by a combination of such means. Foreseeably, power may be even provided by rechargeable means such as solar cells or other external or replaceable internal sources such as batteries. The subscriber transaction terminal equipment is preferably battery powered. All signal distribution apparatus 20, 22, 24, and 26 include a tamper-resistant housing or otherwise secured enclosure as described by U.S. Pat. No. 4,963,966 or secured in a locked equipment closet of an apartment complex or both. If located in a place exposed to the elements, the housing should be water-tight. Also, the housing should be designed to prevent radio frequency leakage. Signal distribution apparatus 20 is uniquely addressable by headend 10 just as is a known converter/decoder. If two bits of a plural bit unique subscriber address are associated with uniquely identifying one plug-in slot for one of four subscriber modules, common control circuitry may be uniquely addressed with remaining address data not used to secure the data communication. Of course, this concept may be utilized to address any number of subscriber modules associated with the signal distribution apparatus. Just as premium programming is transmitted in the clear and since no data communication is necessarily required with a subscriber premises, a subscriber address need not be transmitted in a secure form. Nevertheless, address security may be desirable so long as converter/decoder or other unique address requisite equipment is provided at a premises. Signal distribution apparatus 20 comprises addressable common control circuitry, an optional plug-in special service module and up to four (or more) plug-in subscriber modules. Upon receipt of subscriber specific premium program, subscriber credit or channel authorization data, the data are stored at memory of common control circuitry of signal distribution apparatus 20. Signal distribution apparatus 20 further comprises a diplexer for providing a forward transmission path which is coupled to automatic gain control circuitry of the common control circuitry. The common control circuitry forwards jamming frequency control data to a subscriber module. Channel interdiction circuitry associated with each subscriber module then selectively jams unauthorized premium programming dropped via a particular drop 66 to a particular subscriber. Consequently, signal distribution apparatus 20 is reasonably compatible with downstream addressable authorization data transmission known in the art. No scrambling of premium channels (and no resulting artifacts) is necessary or desirable. Furthermore, no additional forms of service security are necessary such as channel encryption, in-band channel or tier verification or other security measures. The would-be service pirate must attempt to remove a particular pseudo-randomly timed jamming signal placed at a varying frequency or seek to tamper with the off-premises signal distribution apparatus 20 or derive a signal from shielded and bonded cables which should likewise be maintained secure from radio frequency leakage. Two way data transmission is provided via a so-called sub-split frequency spectrum comprising the band 5-30 megahertz for upstream, reverse path transmission toward headend 10 and a spectrum from 54-550 megahertz for downstream forward transmission. In particular, an amplitude shift keyed data transmission signal at approximately 5 MHz is used for communication on drop 66, while a binary phase shift keyed signal is used for upstream data transmission in the T8 band to headend 10. Distribution amplifiers 34 distributed along the distribution plant according to known prior art design techniques separate and separately amplify the two transmission bands. They are distributed along the transmission path in a manner so as to preclude the carrier-to-noise ratio of either transmission path from being two low. Also, at a headend 10, there is located a radio frequency data receiver and data processor for receiving data transmissions from the off- or on-premises subscriber equipment. Details of this equipment are more particularly provided by application Ser. No. 07/498,084, incorporated as necessary by reference. With reference to FIGS. 2 and 3, the signal distribution apparatus 20 of the present invention will be described in further detail. While the method and apparatus according to the present invention will be described in the context of being located at the entry to a subscriber's premises, it will be appreciated that the present invention may be utilized anywhere along the signal path to a subscriber premises where a service pirate may obtain access to the distributed signals. The invention thus may be utilized not only in connection with apparatus 20 but apparatus 34, 44, 22, 24, 46, 26 and so on along the signal path via serving cable or drop to the entrance into the subscriber's premises. The signal distribution apparatus 20 includes an enclosure 200 having a casing 210 and a cover 220. Cover 220 rotates about hinges 231 so as to protect enclosed circuitry within casing 210 from damage. Flange 230 and flange 235 extend around the periphery of casing 210 and cover 220, respectively. Further details of the enclosure can be found in U.S. Pat. No. 4,963,966 and will not be provided here. The enclosure houses circuitry, for example, mother board 240, further described in connection with FIG. 5 below, individual subscriber modules 250, further described below in connection with FIG. 6, and power supply circuitry 246. Mother board 240 and individual subscriber modules 250 are further described below in connection with FIG. 5 and 6. FIGS. 2 and 3 depict the enclosure with a tamper override module (TOM) 400 according to the present invention inserted into interface 245, which may be the same interface to motherboard 240 used for receiving a special service module (not shown). The distribution apparatus is equipped with a plunger 237 which cooperates with boss 236 such that when cover 220 is opened from casing 210 as shown, the plunger, which is conveniently spring mounted, will rise and so open (or close) a contact of an associated tamper detector (FIG. 5). As depicted in FIGS. 2 and 3, the plunger 237 is attached to power supply 246 and the boss 236 is attached to cover 220. However, the invention is not limited in this respect. Plunger 237 and boss 236 may be positioned elsewhere within enclosure 200 to realize the desired functional relationship as would be apparent to one of skill in the art. If service is required and according to the present invention, a service person upon opening the cover 220 inserts TOM 400 (which in alternative embodiments may comprise a magnetic card, a punch card, or other equivalent device) into interface 245. Thereby, a controller may be satisfied that the entry into the enclosure 200 is authorized provided a code transmitted by the TOM 400 is accepted. In a preferred embodiment, TOM 400 is equipped with ears 238 and 239 which overlap a volume immediately above flange 230 after it is inserted into casing 210. In this manner, the service person, should he attempt to close the cover 220, will be precluded from closing the cover and so will be reminded to remove the TOM 400. According to one preferred embodiment, TOM 400 comprises a body portion capable of holding the circuit components described below in connection with FIG. 7. This preferred embodiment of TOM 400 is depicted in FIG. 4. Body portion 410 preferably comprises a card including first and second lateral sides 411, 412 and top and bottom portions 413, 414. The body portion 410 is preferably formed of circuit board material. However, the invention is not limited in this respect. As already suggested, the TOM may simply be a magnetic card such as a plastic credit card, a punched paper card or other equivalent structure sufficing as a key for coupling with a suitably provided interface 245. Lateral sides of TOM 400 are spaced apart a distance equal to the width of casing 210 so as to snugly fit therein as shown in FIGS. 2 and 3. Ears 238, 239 extend from lateral sides 411 and 412, respectively. The ears 238, 239 are preferably trapezoidal as shown in FIG. 4 but may be any shape so as to preclude closure of cover 220. Ears 238, 239 extends over the flange 230 of casing 210 when the TOM 400 is inserted into the enclosure to thereby prevent the cover 220 from closing. Bottom portion 414 includes an extension 417 for insertion into interface 245 coupled to a controller of motherboard 240, for example, microprocessor 560 (FIG. 5). The extension 417 has an electrical connector comprising plural conductive leads 420 printed thereon for coupling with interface 245 to mother board 240 of FIGS. 2 and 3. The conductive leads 420 extend on body portion 410 to the various circuit components, including converter 430, shift register 440, LED 450, and switch 460 as per FIG. 7. Top portion 413 includes interface 480 extending therefrom, for example, for facilitating communications with the mother board by service personnel for diagnostic or other purposes or with the headend 10 or with subscriber terminal equipment. Referring to FIG. 4, extension 417 is slightly off-center off bottom portion 414 and interface 245 is suitably mounted on motherboard 240 to receive TOM 400 in only one direction. Furthermore, components such as interface 480 are mounted in such a way as to preclude installation of TOM 400 in any manner other than that shown in FIG. 2. The common control circuitry of signal distribution apparatus 20 will now be described with reference to the block diagram FIG. 5 for serving four subscriber modules in accordance with the block diagram FIG. 6 and a tamper override module according to FIG. 7. Referring particularly to FIG. 5, feeder cable 28 is shown entering off-premises interdiction apparatus 20 at FEEDER IN and leaving at FEEDER OUT. Power PWR may be provided via the feeder cable, by means of the subscriber drop or locally by internal or external means. Depending on the source of power PWR, input power may be of alternating or direct current. A directional coupler 510 which may be in the form of a plug-in module taps into the broadband serving cable 28. A broadband of radio frequency signals is thus output to highpass filter 520 of diplex filter 595. Highpass filter 520 passes a downstream band of frequencies, for example, 54-550 MHz comprising at least the cable television spectrum and any separate data carrier frequency, such as 108.2 MHz, and blocks the upstream band of frequencies, for example, 5-30 MHz (in a bi-directional application). For an off-premises signal distribution system, the cable television spectrum may particularly comprise a narrower frequency band from about 54 MHz to 350 MHz. Lowpass or bandpass filter 421 passes at least the 0-30 MHz spectrum and more particularly a pass band comprising the T8 band from approximately 14-18 MHz. One of twenty-three data channels may be selected for upstream data transmission from within the T8 band to avoid noisy regions of the spectrum. Circuitry associated with broadband signal "seizure" from the distribution cable 28 may be conveniently mounted on a single board, conveniently named a seizure board of interdiction apparatus 20, more particularly described in FIG. 15 of U.S. Pat. No. 4,963,966, but described in general terms herein as at least comprising the directional coupler 510 and diplex filter 595 of FIG. 5. A common automatic gain control circuit as disclosed in FIG. 5 comprises variable attenuator 530, RF amplifier 533, directional coupler 532, and AGC control circuit 531. This automatic gain control circuit appropriately regulates the broadband RF signal power to fall within established limits. The common circuitry of FIG. 5 is collocated or closely located to the subscriber modules which will be further described in connection with FIG. 6 and may be contained in the same housing with the special service units for each subscriber as described in U.S. patent application Ser. Nos. 07/618,745 and 07/625,901 filed Nov. 27, 1990. Also connected to directional coupler 532 is a data receiver 540 for receiving downstream forward data transmissions from the addressable data transmitter 14 located at headend 10. Data receiver 540 receives data transmitted, for example, over a data carrier of 108.2 megahertz and provides unprocessed data to data decoder 550. In accordance with an established protocol and as briefly described above, such data may be in the form of an operation code (command), a subscriber unique address and associated data. Data decoder 550 processes the data and provides the separately transmitted data to microprocessor 560 for further interpretation in accordance with a built-in algorithm. Microprocessor 560 is most efficiently chosen to alleviate as many responsibilities from any microprocessor provided for an individual subscriber module and so is most conveniently an eight bit microprocessor having eight kilobytes of internal code such as a motorola 68HC05C8 or other suitable microprocessor having internal random access memory and program memory. Received data may be stored in nonvolatile memory (NVM) 570 by microprocessor 560. NVM 570 is preferably three NMC93C46N (64×16) or one NMC93C66N (256×16) by National Semiconductor. However, the invention is not limited in this respect. NVM 570 may store an authorized access code received from the headend, enable/disable information for the tamper system, and/or a tamper detected flag. The tamper detected flag will be stored as a result of operation of tamper detector 555 which provides a contact closure or opening signal over a dedicated lead to microprocessor 560 as a result of the operation of plunger 237 as cover 220 (FIG. 2) is opened. Data may be stored in NVM 570 and jamming frequency control data downloaded when needed to a subscriber module according to FIG. 6 via a serial peripheral interface but 590 connecting microprocessor 560 with separate microprocessors 600 associated with each provided subscriber module as shown in FIG. 6. Furthermore, in an alternative embodiment, separate microprocessors 600 may be replaced by an application specific integrated circuit which performs functions under the control of microprocessor 560 as taught in U.S. application Ser. No. 07/896,292, entitled "Interdiction Method and Apparatus with Programmable Jamming Effectiveness," Jun. 10, 1992. An exemplary application specific integrated circuit comprises a Scientific-Atlanta part No. 463563/463564 manufactured by AMI/Gould of Pocatello, Idaho. Microprocessor 560 is further coupled to interface 245 for communicating with TOM 400 or a special service module, if provided. Furthermore, microprocessor 560 may communicate, for example, upstream frequency and amplitude control data to microprocessors associated with each special service module, as described in U.S. application Ser. Nos. 07/618,745 and 07/625,901. The special service module may share interface 245 with the TOM 400 or, in an alternative embodiment, may be connected to a separate interface. Variable attenuator 530 regulates the received broadband of picture carriers to a reference level while the microprocessor 560 controls the jamming carrier level outputs of associated subscriber units within the prescribed range. Microprocessor 560 consequently interprets both global communications addressed to common control circuitry or communications addressed to unique subscribers for operation of subscriber modules such as service credit, authorization commands, operation commands, or any combination thereof. If appropriate, microprocessor 560 ignores global, group addressed, or specifically addressed communications to other signal distribution apparatus or to conventional converter/decoders. An example of global communications peculiar to signal distribution apparatus 20 is premium channel frequency data for each premium channel or channel over which premium programming at a particular point in time is provided via headend 10. Examples of addressed communications to common control circuitry include communications comprising premium channel or programming authorization information, communications instructing the common control circuitry to provide credit to a particular subscriber, or communications changing the authorized access code of the signal distribution apparatus. Serial peripheral interface buses 590,245 may be two way communications links by way of which microprocessors 600 (FIG. 6) may, at least, provide status reports to microprocessor 560 upon inquiry or TOM (FIG. 7) may supply an access code or request status reports from microprocessor 560. Alternatively, a microprocessor of FIG. 6 may tap into a parallel contention-type bus 590 and bid for communication to either a microprocessor 560 of common equipment or another microprocessor 600 or may directly communicate with any of the other associated microprocessor over a separate serial bus 590. In a similar manner, both a special service module and a TOM may share a bus 245 and likewise bid for communication for access to microprocessor 560 or another microprocessor 600. Radio frequency splitter 580 provides broadband radio frequency signals comprising a broadband cable television service spectrum separately to each subscriber module according to FIG. 6 that is provided. FIG. 6 is an overall block schematic diagram of a subscriber module of signal distribution apparatus 20 including a diplex filter 695. A microprocessor 600 is associated with a particular subscriber module and communicates with microprocessor 560 of FIG. 5 over a serial peripheral interface bus. Microprocessor 600 may comprise an eight bit microprocessor equipped with only two kilobytes of code, this microprocessor being relieved of overall control responsibilities by microprocessor 560. Consequently, microprocessor 600 may conveniently comprise a Motorola 68HC05C3 microprocessor or similar unit. In the alternative embodiment including an application specific integrated circuit, microprocessor 560 may assume the control tasks of microprocessor 600 so that microprocessor 600 can be replaced. Furthermore, analog multiplexer (MUX) 630, D/A converter 620, RAMs, buffers, and prescalars 675 may all be replaced by the application specific integrated circuit. A 5-30 MHz or other lowpass band, more particularly, a 0-15 MHz lowpass band, may be provided for upstream, reverse transmissions from corresponding subscriber equipment on the subscriber premises. Such a reverse path is completed to the subscriber via terminal OS. Also, power may be transmitted up the subscriber drop to the subscriber module of FIG. 6 and withdrawn at terminal OS. The broadband radio frequency television spectrum signal from FIG. 5 is provided to terminal IS. Referring to the path connecting terminal IS to terminal OS, there are connected in series a service denying switch 689, a radio frequency amplifier 687, a jamming signal combiner 685, and a high pass filter 691. Service denying switch 689 is under control of microprocessor 600. In the event of an addressed communication from headend 10 indicating, for example, that a subscriber is to be denied service for non-payment of a bill, service denying switch 689 may be opened thereby disconnecting service. The switch 689 may be closed and opened during scheduled periods of authorized periodic service. In addition or in the alternative, a high frequency amplifier 687 may be powered down under control of microprocessor 600 whenever service is to be denied. Otherwise or even in addition, amplifier 687 may be set at discrete gain levels, under microprocessor control, to provide supplemental gain to be broadband television signal if a subscriber has a plurality of television receivers (TV's and VCR's) over and above a nominal amount. Amplifier 687 may comprise adjustable bandpass filter circuits under microprocessor control for selectably limiting service bandwidth to a subscriber. Alternatively, all subscriber jamming signal generating equipment, for example signal generating equipment 641-644, may be replaced by a programmable bandpass filter 691 which under the control of the microprocessor selectably limits service bandwidth to a subscriber. An appropriate control signal waveform output SDPS is provided by microprocessor 600 for controlling switch 689. Also the same ON/OFF control signal that is used to control the switch 689 may control the powering up and down of amplifier 687 as control signal SDHP. The status of the connect/disconnect condition of switch 689 is preserved in NVM 570 or other memory associated with microprocessor 560 or microprocessor 600. Furthermore, the intended state of switch 689 or related denial circuits is to be always closed or connected. However, the normal state of the condition of switch 689 of a periodic service subscriber is open or disconnected. Continuing the discussion of FIG. 6, jamming signals are interdicted at directional combiner 685 under microprocessor control. Because of the directional characteristic of radio frequency amplifier 687, jamming signals cannot inadvertently reach the common control circuitry of FIG. 5 or the serving cable. Highpass filter 691 of diplex filter 695 prevents any return path signals from reaching combiner 685 and passes the broadband spectrum including any jamming signals toward terminal OS. Reverse path signals, for example, in this embodiment may be radio frequency signals below 30 megahertz. The broadband television spectrum is presumed to be in the 50-550 megahertz range. However, interdiction of premium channel viewing may be allocated anywhere desired within a broader or discontinuous cable television spectrum to be jammed. Consequently, filters 691 and 692 are designed in accordance with this or similarly selected design criteria to block or pass broadband television or reverse path signals as required. Microprocessor 600, responsive to common microprocessor 560, controls the frequency and power level outputs of four (or five if necessary) voltage controlled oscillators 641-644, each of which oscillators jams premium channel frequencies within an allocated continuous range of frequencies. The frequency of the oscillators is set over leads FREQ1-4 in a manner described in U.S. Pat. No. 4,912,760. A power level and ON/OFF operation of the oscillators 641-644 are controlled over leads OPWR1-4. Since premium programming may be transmitted anywhere in the cable television spectrum, the sum of all such allocated portions comprises the entire television spectrum to be jammed (even where non-premium channels are normally transmitted). Also, in accordance with the depicted interdiction system, the television spectrum to be jammed may comprise discontinuous portions or intentionally overlapping portions. A further detailed discussion of frequency control and the interdiction system of FIGS. 1, 5, and 6 may be found in U.S. application Ser. No. 5,014,309. FIG. 7 is a block circuit diagram of the TOM 400. TOM 400 includes an N terminal switch 460 and a voltage source V. The voltage source V is connected via N lines having N resistors, R 1 -R N to switch 460. In one preferred embodiment, N may range from 2 to 8. Switch 460 is also connected to ground G. The N lines are additionally connected to N input of shift register 440. If the number of lines N is less than the number of input pins of the shift register, the remaining input pins are grounded. Switch 460 can be actuated to connect any of the N lines to ground G. Thus, each of the N lines supplies a voltage of either O or V to shift register 440 depending on whether the line is connected to ground G via switch 460. Shift register 440 includes a clock pin and a output pin connected to the microprocessor 560 via lines 710 and 720, respectively. A voltage V is supplied to shift register 440 to supply power thereto. The shift/load pin is connected to converter 430, preferably an RS-232 converter. Line 730 connects microprocessor 560 to two separate pins of converter 430. The first pin receives the signal from line 730, inverts the signal, and outputs the inverted signal to the shift register. The second pin receives the signal from line 730, inverts the signal, and outputs the inverted signal via line 760. Line 740 additionally connects microprocessor 560 to converter 430. Together, lines 710, 720, 730, 740 comprise certain of the conductive leads 420 of extension 417 (FIG. 4). Line 750 connects the converter to, for example, a diagnostic or communications device (not shown) via interface 480 (FIG. 4). Line 760 connects converter 430 to a fist terminal of indicator 450 (an LED as shown here) via resistor 470. A second terminal of indicator 450 is connected to ground G. Indicator 450 may be any device capable conveying information, for example a buzzer, a lamp, liquid crystal display, or the like. The power for the TOM 400 is preferably supplied from the enclosure 200 through power conversion circuit (not shown), for example, via certain of the conductive leads 420. Alternatively, TOM 400 may powered by a battery or in another manner known to those of skill in the art. The preferred operation of the present invention is described in connection with FIG. 8. When cover 220 of enclosure 200 is opened, tamper detector 555 signals microprocessor 560 at step 800. In response to the signal from tamper detector 555, microprocessor 560 determines if tamper mode is enabled from the enable/disable information stored in NVM 570 at step 805. If the tamper mode is disabled, further processes terminate at done box 810. If the tamper mode is enabled, microprocessor 560 sets a cover open bit flag in NVM 570 at step 815. Alternatively, tamper may be permanently enabled thereby obviating steps 805 and 810. Whenever tamper is enabled, it is assumed that the apparatus 20 has entered a tamper diagnostic mode of operation which is at a higher level of priority than any other diagnostic mode the apparatus may enter. At step 820, initializing operations are performed. Depending on the diagnostic mode status of apparatus 20, the LED 450 is powered on if the only diagnostic mode in effect is a tamper diagnostic mode. However, if there is another diagnostic mode in effect other than a tamper diagnostic mode, which is the highest priority diagnostic mode, then, the LED 450 is signaled to flash in accordance with this diagnostic mode. The LED 450 is caused to flash at different rates depending on the diagnostic mode in effect. Other diagnostic modes besides a tamper diagnostic mode may include a refresh timer, a home power counter, an interdiction oscillator malfunction or other diagnostic mode at different levels of priority, the tamper mode being the highest of these levels. Features of apparatus 20 and operation and actuation of the several diagnostic modes are disclosed in copending application Ser. No. 07/896,628, filed concurrently herewith and incorporated herein by reference. A timer internal of the microprocessor 560 is set to a predetermined time for timing a duration of 60 seconds at step 825. The predetermined time is preferably about 60 seconds but may be shorter or longer depending on the typical delay a service person may experience in opening cover 220 and inserting TOM 400. However, the predetermined time may be adjusted by interaction with headend 10 via addressable data transmitter 14. The system operator thus may vary the time interval via addressable data command to controller 560 as experience of the servicemen warrants its variation. Another consideration is that one would not want the duration to be so long as to permit a service pirate to steal service without actuating the tamper diagnostic mode. At step 830 a read TOM code sequence is initiated. That is, the microprocessor 560 outputs a signal set which causes the TOM 400 to output an access code which it then senses. Specifically, the microprocessor 560 outputs a signal on line 730 to the converter 430 to cause enablement of the load data function of shift register 440. In addition, microprocessor 560 clocks the shift register 440 via line 710 to load the data on its input lines. Next, the microprocessor 560 signals the shift register via line 730 and converter 430 to enable the shift function of the shift register. Successive clock pulses applied to the shift register via line 710 causes that data to be serially output on line 720. The microprocessor senses that signal on line 720 to determine if a TOM 400 is installed at step 835. In one embodiment, the read sequence may be repeated twice or even several times in order to ensure that the result is correct. If a TOM 400 is not installed, the microprocessor 560 checks the timer to determine if the predetermined time (here 60 seconds) has elapsed at step 840. If not, the microprocessor 560 determines if the enclosure cover is still open at step 845 by checking the status of the tamper detector 555. If the cover is still open, the process returns to step 830 and the TOM code is again read. If, at step 845, the cover is closed, the microprocessor sets a tamper detected flag in NVM at step 850 and the process proceeds to step 855 where the tamper diagnostic mode is entered. In a preferred embodiment, each subscriber module of the apparatus 20 goes immediately into a tamper diagnostic mode thereby interfering with the service provided to each subscriber connected to the apparatus 20. Interference with the service signal may be accomplished by control of the switch 689. Alternatively, amplifier 687, oscillators 641-644, or other system components may be controlled to interfere with the service provided to the subscribers. The interference may comprise a cutoff of the signal, the supply of a scrambled signal or a signal which is periodically scrambled, the supply of a pulsing signal, or a combination thereof. If the microprocessor 560 determines that the predetermined time has elapsed at step 840, the tamper detected flag is set and the diagnostic mode is entered at steps 850 and 855. The sequence of steps 850 and 855 as well as the sequence of steps 820, 825 may be reversed without adversely affecting the general operation of the present invention. If the microprocessor 560 determines that a TOM 400 is installed at step 835, further action is suspended for 8 to 10 seconds to allow for debounce at step 860. The TOM code is again read at step 865. The microprocessor 560 then compares the access code supplied by the TOM 400 to the authorized access code stored in NVM 570 at step 870. If the codes do not match, the process proceeds to steps 850 and 855 wherein the tamper detected flag is set and the diagnostic mode is entered. If the access code supplied by the TOM 400 matches the authorized access code, the process moves to step 875 and the microprocessor 560 signals the converter 430 via line 730. Responsive to such signal, the converter 430 actuates the indicator 450 (here, the LED is illuminated) to indicate that the access code was correctly entered. If the apparatus is in a diagnostic mode other than a tamper mode, LED 450 is caused to continue flashing the LED at a rate consistent with that diagnostic mode. If, for example, the refresh timer diagnostic mode is in effect and the signal distribution apparatus is in such diagnostics mode, the microprocessor 560 signals the indicator 450 via line 730 and converter 430 to indicate such diagnostic mode. After the diagnostic mode is entered at step 855 or after step 875, the microprocessor 560 waits for a signal from the temper detector 555 indicating that the cover is closed. Thus, according to the preferred embodiment, TOM 400 can be removed to provide a service person with access to motherboard 240 below or for other maintenance. When the cover closed signal is received, the microprocessor 560 rests the cover open flag in NVM 570 at step 885 and the process is completed at done step 890. In an alternative embodiment, it is suggested that the tamper detector 555 and tamper plunger 237 may be operated as a signaling mechanism in the absence of TOM 400. In this embodiment of apparatus 20, a predetermined code may be signaled as if plunger 237 were a pulsing device. Microprocessor 560 then may operate to determine a several digit code which must be input in accordance with predetermined criteria. For example, the code must be tapped into the apparatus 20 at a particular predetermined rate for each digit with a predetermined delay that must be exceeded between entry of each digit. If the code is accepted, the apparatus 20 may be satisfied and a tamper diagnostic mode avoided. Alternative embodiments and operation of the tamper override capability are now discussed. For example, the tamper status flag indicating an open cover may be re-enabled 60 seconds after TOM 400 is removed from interface 245 rather than after cover 220 is closed. Of course, re-enablement of the tamper status flag may be triggered according to the occurrence of another event as will be apparent one of skill in the art. In any of the above described embodiments involving a plunger 237, the plunger 237 may be provided with a cap or a cover to prevent accidental actuation by the service person whenever the service person is working on apparatus 20. While the principles of the present invention have been described above in conjunction with specific apparatus, it is to be understood that this description is made only by way of example and not as a limitation to the scope of the invention.

Apparatus for broadband signal distribution is equipped with a detector for detecting when the apparatus is open and is further adapted to receive an entry code generator. If the entry code generator is not installed or, if installed, fails to generate an entry code which matches an authorized entry code stored in the apparatus, the apparatus enters a particular mode of operation whereby service provided by the apparatus is interfered with. The entry code generator comprises a circuit board for coupling with an electrical connector of the apparatus having at least one ear which precludes closure of a cover to the apparatus. A method of screening access to the apparatus from others than authorized service personnel includes waiting a predetermined period of time for coupling the access code generator before service is interfered with. Furthermore, the stored access code may be changed or the predetermined time varied by addressable command from a headend of a broadband signal distribution system.

CLAIM FOR PRIORITY [0001] This application claims the benefit of priority to international application No. PCT/DE2003/04130, which was published in the German language on Dec. 10, 2003, which claims the benefit of priority to German Application No. 103 05 415.4, filed on Feb. 6, 2003, the contents of which are hereby incorporated by reference. TECHNICAL FIELD OF THE INVENTION [0002] The invention relates to a method and an apparatus for media-redundant operation of a terminal in a network. BACKGROUND OF THE INVENTION [0003] The expression network component covers any apparatus which can send and/or receive data packets in the network. In this sense, network components may, for example, be switches or hubs. A switch is an active network component for connecting individual segments of a network to one another. A hub is likewise an active network component, for example a star coupler. The expression connecting medium is used for components of a network via which data and signals can be physically transmitted between different network components. These include, for example, twisted pair cables, coaxial cables and fiber-optical cables. The expression medium covers not only the pure connecting medium but also its connecting components to the network on the one hand (network components) and to the terminal on the other hand (transmission devices). However, since the connecting media in particular may be affected by possible failures, for example as a result of external destruction, these represent the major focus for the design of a media-redundant connection. If a terminal is connected to the network via two physically independent media (media redundancy), then, if one of the media fails, safe and reliable communication is nevertheless possible with the terminal. Since disturbances to individual media are unavoidable, manufacturing failures can be avoided by means of media redundancy in the case of media disturbances, for example in production operations. [0004] Media redundancy can be achieved in a terminal by, for example, equipping the appliance with two transmission devices each having their own media access controller, which controllers are connected to the network via different connecting media. This configuration can be provided in the terminal, for example in the form of two separate network boards. In this case, conventional network components may be used for the inclusion of a terminal. Furthermore, failure monitoring is possible even into the terminal itself, from application to application, in all network protocol layers separately for the two different network links. However, this has the disadvantage that two network boards are necessarily required in each terminal, and this leads to an increase in the costs of the terminals. Furthermore, in this situation, each terminal must be assigned two media access control addresses (MAC addresses) as well as two Internet Protocol addresses (IP addresses). Failure monitoring and a change between the two connections via the two network boards takes place at a higher network protocol layer, generally at the application layer, and is thus done using computer capacity in the terminal. Failure monitoring at the application layer thus results in the computer power of the terminal being restricted with regard to the actual application. The switching times in the event of a disturbance in one of the media are typically in the region of 100 ms. [0005] Another implemented approach for achieving media redundancy is referred to as the ring redundancy approach. This approach is described, by way of example, in the document “Industrial Ethernet™ startet durch—Switching und 100 Mbit/s in der industriellen Kommunikation” [Industrial Ethernet™ starts by—switching and 100 Mbit/s in industrial communication], K. Glas, 1998. In this case, one or more terminals is or are connected to a ring network. If this network is interrupted at one point, then all of the terminals and network components in the ring network are still connected to one another, such that they are still fully serviceable. In this solution, the terminals each require only one network board with one transmission device and one media access controller. In the event of a disturbance, switching takes place at a so-called second network protocol layer (“layer 2 ”). At the application layer, no program is required for failure monitoring, although no monitoring from application to application is possible, either. Another disadvantage is that the switching times are in the region of 200 ms. A further disadvantage is that it is not possible to use standard switch components to produce the network. The only switch components which may be used are those which are specially designed for use in a network with ring redundancy. SUMMARY OF THE INVENTION [0006] The invention relates to a method and an apparatus for media-redundant operation of a terminal in a network. When electronic terminals are being linked, for example computers or electronically controlled manufacturing machines, by means of a network, it is desirable for media redundancy to be provided for each terminal. A terminal may in this case be any desired device which has a network component of the network or is connected to such a component. [0007] In one embodiment of the invention, there is a method and apparatus for safe, reliable, low-cost and fast media-redundant operation of a terminal in a network, on the basis of media redundancy achieved by duplication of the transmission devices and of the associated media access controllers. [0008] In another embodiment of the invention, there is a terminal connected to a network via two transmission devices in a physical bit transmission layer, of which devices one is active at a time, in order to interchange data with the network. The transmission devices are connected to a single control device. If the first (active) of the two transmission devices finds that there is no good connection between it and the network via a first connecting medium, it is deactivated by the control device, the further transmission device is activated as the active transmission device, and an electronic failure message is sent into the network, so that the network can adapt itself to the change that has occurred with respect to the activity status of the two transmission devices, and its consequences with regard to the passing on of data for that terminal. The connection between the terminal and the network is then provided by means of the second (redundant) connection via the further transmission device and the further connecting medium. [0009] One advantage of the invention is that a method and an apparatus are provided which allow very fast switching of the redundant link between a terminal and a network. The expression “fast switching” describes the fact that the data interchange between the terminal and the network, which is handled via a first connecting medium, can be continued after a very short time, that is to say quickly, via the further connecting medium in the event of a disturbance or an interruption in this first connecting medium. Since the change from the first connecting medium to the further connecting medium requires only the deactivation of the first transmission device, for which the disturbance was found, and the activation of the further transmission device in the physical bit transmission layer, as well as the sending of an electronic failure message, the change can be carried out in less than about 1 ms. The switching need not be initiated at a higher network protocol layer. The invention may not only be implemented in an electronic circuit, but may also be implemented by means of software, which interacts with an electronic circuit. An apparatus according to the invention can be produced at low cost and is suitable not only for use in appliances to be designed from new, but also for use with old appliances, in order to connect them to the network with media redundancy. [0010] Still another embodiment of the invention comprises a data packet to be transmitted from the terminal to the network being analyzed by an analysis device, a media access control address being determined in the process and being assigned to a media access controller for the terminal, and the media access control address being inserted into the electronic failure message during the production of the electronic failure message. These method steps make it possible to carry out the method independently of any prior knowledge of the media access control address of the media access controller for the terminal. [0011] One advantageous embodiment of the method according to the invention furthermore provides for the network to analyze the failure message and to adapt itself to it in such a way that data packets which are addressed to the terminal are passed on to the further connecting medium. This means that the network can adapt itself in a very short time to the change in the connecting medium via which the terminal is connected to the network. This avoids unnecessary network traffic, and data packets are no longer passed to the connecting medium which is faulty or is subject to a disturbance. [0012] Another embodiment of the method according to the invention provides that the network component analyzes the failure message and adapts its address addressing table for the media access controller for the terminal in such a way that data packets which are addressed to the terminal are passed on to the further connecting medium. This allows particularly fast switching of the network, since only the settings of the network component need be changed. [0013] Furthermore, one advantageous embodiment of the method according to the invention comprises the use of a network component, to which the first and the further connecting medium are connected. In this case, the network component just has to internally redirect the data packets that are addressed to that terminal to another connection in order to switch the connection to that terminal from the first (disturbed) connecting medium to the further connecting medium. This can be done particularly easily and quickly when using a network component that is designed appropriately for this purpose. [0014] However, one particularly advantageous embodiment for the method according to the invention is regarded as the use of a further network component, to which the further connecting medium is connected. This also achieves redundancy with regard to a failure of one of the network components to which the terminal is connected directly. [0015] A further embodiment of the method consists in that a communication is carried out from the terminal with the network component by means of a standard network protocol, in particular an Internet Control Message Protocol (ICMP), in order to check the serviceability of the network component. This allows the terminal additionally to react even to those failures of the network component in which a sound physical connection still exists between the transmission device and the network component, despite the failure. [0016] It is also advantageously possible to provide for an application to be run on the terminal, and for the application to interchange checking data packets after predetermined time intervals with another network component which is connected to that network component, in order to determine whether the network component is serviceable. This means that it is possible to monitor the failure of the network component when the latter is not able to communicate by means of the standard network protocol. [0017] Another embodiment of the method provides for the other network component to produce an application failure signal if it has not received any of the checking data packets after a predetermined number of predetermined time intervals. This makes it possible to detect a failure of the application on the other network component. If, for example, the application in the terminal is instructed to interchange data with the other application in the other network component, then it is important to be able to distinguish between “network disturbances” and application failures, in the event of communication disturbances. [0018] With regard to the apparatus according to the invention, one advantageous embodiment provides that the first and the further connecting medium are connected to the network component. This makes it possible to switch particularly quickly to the second connecting medium in the event of a failure of the first connecting medium, since the network component just has to switch internally to another connection. However, an appropriately designed network component must be used for an embodiment such as this. [0019] One alternative advantageous embodiment of the apparatus according to the invention provides that the further transmission device is connected to another network component via the further connecting medium. This means that the apparatus is also provided with redundancy with regard to a failure of the network component or of the other network component. [0020] Another advantageous embodiment of the apparatus provides that the transmission device, the further transmission device and the control device are arranged on a network board in the terminal. This allows the apparatus to be integrated in a terminal in a simple manner. This also makes it possible to connect commercially available personal computers with media redundancy to a network. [0021] One embodiment consists in that the transmission device, the further transmission device and the control device are arranged in a redundancy switch, which is connected to the terminal. A redundancy switch such as this allows terminals which are not designed for this purpose to be connected with media redundancy, without having to modify the terminals. [0022] A further embodiment of the apparatus provides for the redundancy switch to have an additional transmission device in the physical bit transmission layer, in order to interchange data packets with the terminal. This development allows a network-compatible terminal to be connected to the network with media redundancy, without any additional precautions. In addition, this development can be used to connect a terminal not only via two separate connecting media, but also via two different connecting media types, for example via a twisted pair cable and via a fiber-optical cable. [0023] A further advantageous embodiment of the apparatus provides an analysis device for analysis of a data packet to be transmitted into the network from the terminal, in order to determine a media access control address and in order to insert the media access control address into the electronic failure message. This embodiment can be operated without any prior knowledge of the media access control address of the terminal, and is particularly advantageous when the media access control address changes as a result of replacement of the media access controller in the course of servicing or upgrading work on the terminal, since the new media access control address is determined automatically in this embodiment. [0024] One embodiment of the invention is for the control device to have the analysis device. This results in a compact embodiment. [0025] Another advantageous embodiment of the invention provides for the control device to be an electronically programmable logic device (ELPD). A control device such as this costs little and can easily be matched to a media access control address for the terminal. This development means that it is possible to change particularly quickly from the connecting medium to the other connecting medium. [0026] Another embodiment of the apparatus provides for the control device to comprise the media access controller. This development makes it possible to design a very compact and low-cost embodiment of the invention. BRIEF DESCRIPTION OF THE DRAWINGS [0027] The invention will be explained in more detail in the following text using exemplary embodiments and with reference to the drawings, in which: [0028] FIG. 1 shows a terminal which is connected with media redundancy to a network. [0029] FIG. 2 shows a schematic illustration of a network. [0030] FIG. 3 shows another embodiment of a terminal, which is connected with media redundancy to a network. [0031] FIG. 4 shows an apparatus for media-redundant operation of a terminal in a network. [0032] FIG. 5 shows a further apparatus for media-redundant operation of a terminal in a network. DETAILED DESCRIPTION OF THE INVENTION [0033] FIG. 1 shows a terminal 1 which is connected with media redundancy to a network 4 . By way of example, it is assumed here that the network 4 is an Ethernet network. However, the invention can also be used in other network types. The terminal 1 has a processor (CPU) 2 , which controls the terminal 1 and runs applications on the terminal 1 . The processor 2 is connected to a control device 5 which, in this embodiment, comprises a media access controller (MAC) 3 . The media access controller 3 controls the connection of the terminal 1 to the network 4 on a so-called second network protocol layer (“layer 2 ”). The media access controller 3 is assigned a unique media access control address (MAC address), by means of which the media access controller 3 for the terminal 1 is distinguished from other media access controllers for further terminals in the network 4 . The control device 5 is connected to a transmission device 6 for a so-called physical bit transmission layer (“physical layer device”-PHY) and to a further transmission device 7 in the physical bit transmission layer. The transmission device 6 and the further transmission device 7 convert the data transmitted to them from the media access controller 3 into physical signals, which can be transmitted via a first connecting medium 8 and via a further connecting medium 10 . The connecting media may comprise coaxial cables, twisted-pair cables, glass-fiber cables, wireless connections, etc. The transmission device 6 , the further transmission device 7 and the media access controller 3 may expediently be combined on one network board 17 . [0034] The transmission device 6 is connected to a network component 9 via the first connecting medium 8 . The further transmission device 7 can likewise be connected to the network component 9 via the further connecting medium 10 . In a situation such as this, media redundancy exists for the connection of the terminal to the network only with regard to the connecting media 8 and 10 . In this situation, an appropriate network component 9 must be used, which internally converts the connection to the terminal to the other connecting medium in the event of a failure of one of the connecting media. [0035] However, the following text will consider the situation as illustrated in FIG. 1 , in which the further transmission device 7 is connected to another network component 11 via the further connecting medium 10 . The network component 9 and the other network component 11 may each in particular be a switch, a HUB or the like. The network component 9 and the further network component 11 are connected to one another via another connecting medium 12 . [0036] During operation, only the first transmission device 6 or the further transmission device 7 is in each case active. The first transmission device 6 and the further transmission device 7 are activated/deactivated by the control device 5 in this embodiment. The first transmission device 6 and the further transmission device 7 carry out autonomous automatic monitoring, in order to determine whether there is a sound physical connection via the appropriate first connecting medium 8 and the further connecting medium 10 , respectively, to the network component 9 and to the other network component 11 . This check is referred to as a connection status check or a link status check. During this check, the MAC address that is assigned to the media access controller 3 is not transmitted to the first network component 9 and to the other network component 11 . [0037] Those skilled in the art will be familiar with the method of operation of the link status check by the first transmission device 6 and the further transmission device 7 in the physical bit transmission layer. By way of example, the link status check will be explained for the first connecting medium 8 . It is assumed that the first connecting medium 8 comprises an optical fiber and a further optical fiber, as is normal for optical connecting media. The first transmission device 6 uses the optical fiber to send test signals based on a protocol for the bit transmission layer to the network component 9 . During normal operation, the network component 9 “reflects” the test signals to the further optical fiber. If it is now also assumed that the first connecting medium 8 is detached from the network component 9 , then the latter identifies from the lack of the “reflected” test signals that the first connecting medium 8 has been disturbed. A so-called “far end fault” is present. [0038] In order to ensure that either the first transmission device 6 or the further transmission device 7 is active at one time, and sends data packets using the MAC address assigned to the media access controller 3 , it is possible for the network 4 to configure itself such that the data packets directed to the terminal 1 from the network 4 are transmitted to the active transmission device of the physical bit transmission layer. [0039] In order to explain the reaction to a media failure, it is assumed that the first transmission device 6 is active initially. If a disturbance now occurs on the first connecting medium 8 , then the first transmission device 6 finds during the link status check that there is no longer a sound physical link to the network component 9 . The failure of the connection is transmitted to the control device 5 , as is indicated by an arrow 13 . Furthermore, the fault can also be transmitted to other components of the terminal, for example the processor and to programs in higher network protocol layers, in particular the application layer, as is indicated by arrow 14 . [0040] As soon as the connecting media failure is found by the first transmission device 6 , it is deactivated by the control device 5 and the further transmission device 7 is activated, as is indicated by other arrows 15 a and 15 b . In order to allow the network 4 to continue to transmit the data packets intended for the terminal 1 via the further connecting medium 10 to the terminal 1 a failure message is first of all sent via the other network component 11 to the network. The failure message is advantageously in the form of a so-called multicast message. This message therefore contains no addressee and is passed on to all of the network components in the network 4 which belong to the network in the physical bit transmission layer. The network can also be formed by a network element which comprises all of those network components to which data packets of the terminal 1 can be passed on exclusively on an MAC address basis. [0041] Intelligent network components in the network, in particular switches and other terminals, can adapt their MAC address addressing tables on the basis of the failure message, which includes the MAC address assigned to the media access controller 3 . The expression intelligent network component in this case covers any network component which can analyze the data packets received by it and is able to react in some way to a result of the analysis. [0042] The analysis of the failure message and the adaptation of the addressing tables can be carried out by any intelligent network components or, by way of example, predominantly by the network component 9 . After a failure of the connection to the terminal 1 via the first connecting medium 8 and the first transmission device, this network component 9 can pass on the data packets intended for the terminal 1 via the other connecting medium 12 to the other network component 11 . [0043] In the situation where the two transmission devices 6 and 7 are connected via the connecting media 8 and 10 to the one network component 9 , the data packets which are intended for the terminal 1 is passed on internally in the network component 9 to a connection which is connected to the further connecting medium 10 . [0044] The deactivation and the switching from the first transmission device 6 to the further transmission device 7 can be carried out in the second network protocol layer (layer 2 ) or in a higher network protocol layer. The switching can also be carried out by a control device 5 , which need not comprise the media access controller 3 . The failure message can likewise be produced in the second network protocol layer (layer 2 ), for example by means of an electronically programmable logic device (“Electronic Programmable Logic Device” (EPLD)) or an application-specific integrated circuit (ASIC), or can be produced in a higher network protocol layer. The failure message can be coded in such a way that the media failure and switching operations are indicated to all the communication partners. The production of the failure message may also include data being read in from a memory, in particular from a read only memory (ROM). [0045] In addition to the monitoring of a failure of one of the connecting media 8 , 10 , it is also possible to monitor for a fault in the network component 9 and in the other network component 11 . If the network component 9 and the other network component 11 are each in the form of an intelligent component, then failure identification can be ensured by means of cyclic sending of requests to the network component 9 and to the other network component 11 on the basis of a standard protocol, such as the ICMP (Internet Control Message Protocol). The correctly operating network component 9 and the other network component 11 respond to the requests. [0046] FIG. 2 shows an enlarged detail from the network 4 . The same features in FIGS. 1 and 2 are provided with the same reference symbols. In the exemplary embodiment shown in FIG. 2 , an additional network component 23 or further terminals 25 , 26 , which are likewise connected directly to the network component 9 and to the other network component 11 , respectively, can be checked in the manner described above. [0047] The check for failure identification via an ICMP communication can selectively also be triggered or initiated by an application in the terminal 1 using the so-called watchdog method (monitoring method). In this case, a communication based on the standard protocol, in particular an ICMP communication, with the additional network component 23 or with one of the further terminals 25 , 26 which are connected directly to the network component 9 to which the transmission device 6 (which is assumed to be active in this case) is connected, is carried out at specific time intervals. The additional network component 23 or one of the further terminals 25 , 26 can use the lack of this communication in accordance with the standard protocol to deduce that the network component 9 has failed. [0048] The application in the terminal 1 can likewise produce an “other component failure signal” if no more request data packets are received by the application in the terminal 1 in accordance with the standard protocol, when such request data packets are initiated by the other application in the additional network component 23 or by one of the further terminals 25 , 26 . [0049] In the embodiment shown in FIG. 1 , the first transmission device 6 and the further transmission device 7 were connected to the network component 9 and to the further network component 11 . FIG. 2 shows an embodiment of an additional terminal 24 , in which the transmission device 28 and the further transmission device 29 are connected to the additional network component 23 via two different connecting media, a connecting medium 40 and a further connecting medium 41 . This embodiment also ensures connecting media redundancy for the additional terminal 24 . In this embodiment, however, the redundancy is lost in the event of a total failure of the additional network component 23 . [0050] In the embodiment shown in FIG. 1 , the network component 9 was connected directly to the other network component 11 via the other connecting medium 12 . FIG. 2 shows a different network topology. The additional network component 23 is arranged between the network component 9 and the other network component 11 , and subdivides the other connecting medium 12 into two parts 12 a , 12 b . Further additional network components can also be arranged between the network component 9 and the other network component 11 , provided that the configuration ensures that multicast data packets from the further additional network components are passed on from the network component 9 to the further network component 11 , and vice-versa. [0051] If the further terminals 25 , 26 are intended to be connected to the network with media redundancy, in addition to the terminal 1 , then it is worthwhile providing a first connecting medium 43 to the network component 9 , and a second connecting medium 44 to the other network component 11 , from each of the further terminals 25 , 26 . This creates a clear network topology, and additionally means that the terminals are also redundantly connected to the network in the event of a failure of the network component 9 or of the other network component 11 . [0052] FIG. 3 shows another embodiment of the invention. Identical features in FIGS. 1 to 3 are provided with identical reference symbols. In this embodiment, a control device 5 is arranged between the first transmission device 6 and the further transmission device 7 on the one hand, and the media access controller 3 on the other hand. The control device 5 may be an electronically programmable logic device (EPLD) or an application-specific integrated circuit (ASIC). In this embodiment, the failure signals from the first transmission device 6 and the further transmission device 7 are transmitted to the control device 5 . In the event of a connecting media failure, the control device 5 carries out the “switching”. This may be done transparently for higher network layers. This allows very fast switching. The switching time is in the region of about 1 ms or less. In addition, in this embodiment as well, the connecting media failure can also be transmitted to the higher network protocol layers. [0053] FIG. 4 shows a further embodiment of the invention. Identical features in FIGS. 1 to 4 are provided with identical reference symbols. In this embodiment, a redundancy switch 20 comprises the first transmission device 6 , the further transmission device 7 and the control device 5 . The redundancy switch 20 is separated from the terminal 1 and is connected to it via a connecting line 21 . A data packet to be sent from the terminal 1 to the network 4 comprises the MAC address which is associated with the media access controller 3 for the terminal 1 . An analysis device 22 for the redundancy switch 20 , which is connected to the control device 5 , determines the MAC address which is required for production of the failure message, on the basis of the data packet to be sent. The control device 5 operates in the manner that has been described above for the embodiment shown in FIG. 1 . The redundancy switch 20 allows terminals which are not designed in a redundant form to be easily connected with media redundancy to the network 4 . [0054] FIG. 5 shows a further embodiment of the invention, similar to that illustrated in FIG. 4 . Identical features in FIGS. 4 and 5 are provided with identical reference symbols. In addition to the embodiment shown in FIG. 4 , the redundancy switch shown in FIG. 5 has an additional transmission device 31 in the physical bit transmission layer, which is connected to the control device 5 . The additional transmission device 31 is connected via an additional connecting medium 33 to another transmission device 32 in the physical bit transmission layer. The other transmission device 32 is formed by the network board 17 . The embodiment shown in FIG. 5 is distinguished in that the redundancy switch 20 is designed such that it can be used with the terminal 1 , which is not designed in a redundant form, in order to connect the terminal 1 with media redundancy to the network 4 . Furthermore, the control device 5 in this embodiment comprises the analysis device 22 , as is indicated by a dashed line.

The invention relates to a method and device for medium-redundant operation of a terminal in a network. The device comprises: a first transmission device ( 6 ) of a physical bit transmission plane which is connected to a network component ( 9 ) via a first connection medium ( 8 ) disposed between the first transmission device ( 6 ) and a network component ( 9 ) and which can produce an error signal during automatic monitoring of link to the network component ( 9 ) if no physical link or if a defective link to the network component ( 9 ) exists; another transmission device ( 7 ) on the physical bit transmission plane, which is connected to the network ( 4 ) via another connection medium ( 10 ); and a control device ( 5 ) which is connected to the first transmission device ( 6 ) and the other transmission device ( 7 ) in order to deactivate the first transmission device ( 6 ) as a reaction to the error signal therefrom, enabling activation of the other transmission device ( 7 ) and enabling an electronic error telegram to be sent to the network ( 4 ) via the other connection medium ( 10 ) by means of the other transmission device ( 7 ). The invention also relates to a method for operating said device.

CROSS-REFERENCE TO RELATED APPLICATIONS This application is a Divisional of the U.S. patent application Ser. No. 12/571,110 filed on Sep. 30, 2009, which claims foreign priority to 097150574 filed on Dec. 24, 2008, all of which is hereby incorporated by reference in its entirety. Although incorporated by reference in its entirety, no arguments or disclaimers made in the parent application apply to this divisional application. Any disclaimer that may have occurred during the prosecution of the above-referenced application(s) is hereby expressly rescinded. Consequently, the Patent Office is asked to review the new set of claims in view of the entire prior art of record and any search that the Office deems appropriate. FIELD OF THE INVENTION This invention relates to a polymer molecular film, a photo-electronic element and a method for manufacturing the same. BACKGROUND OF THE INVENTION Polymer is constructed by repetitive chemical bonding between low molar mass molecules. The common structures can be linear, network, or branched. Although conventional polymeric materials are generally insulators, polymers with conjugated chain structures are capable to conduct electricity by transport of pi electrons. Particularly, with the procedures of redox reactions, the conductivity of conjugated polymers can reach the levels of doped inorganic semiconductors or some conductors. This unique property, when combined with other important advantages such as low material cost, simple fabrication processes, compatibility with large area manufacturing, light weights, and bendable mechanical properties, has made conjugated polymers emerging as the vital candidate for next generation optoelectronics. For example, PLEDs (polymer light emitting diode) are the application that is widely studied. In short, conjugated polymers are used as a kind of light emitting material, which in turn is applied between positive electrodes and negative electrodes to form light emitting films. When forward bias is applied, holes are injected into the polymer molecular film from the positive electrodes and enter valence band to become positive polarons. Moreover, electrons are injected from negative electrodes and enter a conducting band to become negative polarons. And the two polarons move in the opposing directions to be combined to emit fluorescence (visible light). Polymer LEDs can become polymer semiconductor laser with proper design and manufacturing process. The operation principle of the polymer semiconductor laser is generally similar to that of polymer LEDs, but that the resonant cavity structure is particularly introduced and population inversion is achieved by suitable electron levels, so that when the light is transmitted in the semiconductor polymer layer, energy gap wavelength photons are stimulated to emit high intensity coherent light. For example, the elements having the similar structure to the aforementioned structures can also be used to generate electrical power, such as using that to manufacture an electrical power generating element using solar energy. Using the energy of the photons to separate the electrons and the holes. After being separated, the holes move towards positive electrodes, and the electrons move towards negative electrodes so as to form charges needed by external circuits, and the photonic energy can be transformed into electrical energy. No matter it's the aforementioned light emitting or electrical power generating elements, better efficiency is needed for increasing the applicability thereof. SUMMARY OF THE INVENTION Therefore, the first objective of this invention is to provide a polymer molecular film. Particularly, the polymer molecular film has higher light-emitting/power-generating efficiency, which can be widely applied in manufacturing all kinds of photo-electronic elements. According to an embodiment, the polymer molecular film is formed on a substrate through a deformation process. Particularly, the polymer molecular films comprise a plurality of conjugated polymers, wherein at least one of the plurality of conjugated polymers has a stretched molecular structure. The second objective of this invention is to provide a method for manufacturing aforementioned polymer molecular films. According an embodiment, the method includes the following steps: First, applying a conjugated polymeric material on a substrate to form a polymeric material layer; then, un-stabilizing the polymeric material layer to form a polymer molecular film; wherein, at least one conjugated polymer of the polymer molecular films has a stretched molecular structure. In an embodiment, the step of un-stabilizing the polymeric material layer to form the polymer molecular film further comprises: disposing the polymeric material layer in vapor of a solvent. In another embodiment, the step of unstablizing the polymeric material layer to form the polymer molecular film further comprises: heating the polymeric material layer until the temperature is higher than glass transition temperature of the conjugated polymeric materials. The third objective of this invention is to provide a photo-electronic element, comprising the aforementioned polymer molecular films, so that the photo-electronic element has higher light-emitting/power-generating efficiency. According to an embodiment, the photo-electronic element comprises a substrate, the polymer molecular film and a protection layer. The substrate and the protection layer can also be the transporting layer for the electrodes, or electrons and holes. The polymer molecular film is formed on the substrate, and the protection layer is formed on the polymer molecular film to prevent the polymer molecular film from oxidation or wear. The fourth objective of this invention is to provide the method for manufacturing the aforementioned photo-electronic elements. According to an embodiment, the steps of the method include the following: First, preparing a substrate, then applying a conjugated polymeric material on the substrate to form a polymeric material layer, and then un-stabilizing the polymeric material layer to form the polymer molecular film; wherein, at least one conjugated polymer in the polymer molecular films has a stretched molecular structure. The advantages and the spirits of this invention can be further understood with the following description and the appended drawings. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 illustrates a flow chart of an exemplary polymer molecular film manufacturing method in accordance with embodiments of this invention. FIG. 2A illustrates an exemplary substrate in accordance with embodiments of this invention. FIG. 2B illustrates an exemplary polymeric material layer in accordance with embodiments of this invention. FIG. 2C illustrates an exemplary polymer molecular film in accordance with embodiments of this invention. FIG. 3 illustrates an exemplary flow chart of the polymer molecular film manufacturing method in accordance with embodiments of this invention. FIG. 4 illustrates an exemplary polymer molecular film in accordance with embodiments of this invention correspond to the photoluminance spectra of different dewetting time during dewetting process. FIG. 5 illustrates the cross-sectional view of an exemplary light emitting element in accordance with embodiments of this invention. DETAILED DESCRIPTION OF THE INVENTION This invention provides a sort of polymer molecular film, photo-electronic element, and the method of manufacturing the polymer molecular film and the photo-electronic element. The followings are the embodiment and the practical applications of the invention, and those will be further described to better explain the characteristics, spirits and advantages of this invention. Please refer to FIG. 1 and FIG. 2A to FIG. 2C . FIG. 1 illustrates the flow chart of the polymer molecular film manufacturing method of an embodiment of this invention. FIG. 2A illustrates the substrate 10 formed by Step S 50 in FIG. 1 ; FIG. 2B illustrates the polymeric material layer 12 formed by Step S 52 in FIG. 1 ; and FIG. 2C illustrates the polymer molecular film 14 formed by Step S 54 in FIG. 1 . As illustrated, the manufacturing method according to this embodiment includes the following steps: Step S 50 , preparing a substrate 10 , for example but not limited to, a glass substrate, Indium Tin Oxides (ITO), etc. Step S 52 , applying conjugated polymeric material to form a polymeric material layer 12 on the substrate 10 . In practice, the conjugated polymeric material is applied to the substrate 10 by spin coating or other proper means. In practice, the thickness of the polymeric material layer 12 is less than 300 nm, such as 200 nm, 100 nm, 50 nm, and preferably less than 30 nm. In practice, the conjugated polymeric material can be a single polymer or several polymers composed of a sort of mixture or copolymer. The conjugated polymer includes, but not limited to, poly[1-methoxy-4-(2′-ethyl-hexyloxy)-2,5-phenylene vinylenel] (MEH-PPV), polythiophene or polyphenylene, etc. Moreover, in practice, the conjugated polymeric material may include other nonconjugated polymeric materials, such as polystyrene (PS), as well as other proper additives. Step S 54 , heating the polymeric material layer 12 such that the temperature of the polymeric material layer 12 is higher than the glass transition temperature (Tg) so as to allow the polymeric material layer 12 to deform to the polymer molecular film 14 . q As illustrated by FIG. 2C , the polymer molecular film 14 has a plurality of protrusion areas 140 and, possibly, a plurality of serial indentation areas 142 . Particularly, at least one conjugated polymer of the plurality of indentation areas 142 has a stretched molecular structure in unilateral direction. Furthermore, in practice, the distance between the surface 1420 of the indentation area 142 and the surface 102 of the substrate 10 is less than 100 nm; for example, between 0.5 nm and 50 nm, and preferably less than 20 nm. In practice, by surface treatment to the substrate surface, the thickness of the indentation areas can be controlled. In practice, Step S 54 is the conventional dewetting step. This step can drag conjugated polymers on the surface of the substrate to form a film having stretched molecular structures, and the originally flat polymeric material layer is now fractured to form droplets (same as the protrusion area described above). Because the conjugated polymers are stretched, the conjugated polymeric chains are less likely to be bent, so that the charges (including electrons and holes) can move freely on the polymer chains without being trapped by polymer chains and therefore, it is easier to emit light/generate power. After the preliminary experiment, the luminous efficiency of the polymer molecular film after being dewetted has at least one order of magnitude higher than that of the polymer molecular film before being dewetted. Please refer to FIG. 3 . FIG. 3 illustrates a flow chart of the polymer molecular film manufacturing method of a different embodiment of this invention. As illustrated by FIG. 3 , Step S 54 in FIG. 1 can be replaced by Step S 56 which proceeds the followings: under room temperature, disposing the polymeric material layer in the vapor of a certain solvent (such as Toluene, p-toluene, Tetrahydrofuran (THF), methanol or other proper solvents), such that the polymeric material layer 12 is un-stabilized and deformed to become the polymer molecular film 14 . Moreover, in practice, the aforementioned Step S 52 and S 54 /S 56 can be repetitively implemented to form multiple polymer molecular film layers to increase light-emitting/power-generating efficiency. Moreover, in practice, the method of this invention can also include the following steps: removing the parts of the protrusion areas which are higher than the indentation areas, such that the polymer molecular film has the surface that is generally flat. By doing this, the situation that the light emitted from the indentation area is affected by the protrusion areas can be avoided. On the contrary, when Step S 52 and S 54 are implemented repetitively to form multiple polymer molecular film layers, the protrusion areas may be remained, so as to help fix the bonding between the films. Please refer to FIG. 4 . FIG. 4 illustrates polymer molecular films of an embodiment of this invention during dewetting process, corresponding to the photoluminance spectra in different dewetting time. In this embodiment, luminous polymeric material MEH-PPV is spin coated on a silicone oxide to form a film with 20 nm thickness, and disposed in an environment of 100 degrees (° C.) of temperature for different time periods to implement dewetting. Moreover, FIG. 4 shows the results generated by a fluorescence spectrometer. As illustrated by the figure, when the dewetting time period is 300 minutes, the film is completely dewetted and the luminous efficiency increases 30 times to that of the film before dewetting. This invention further provides method for manufacturing the photo-electronic element having therein the aforementioned polymer molecular films. The photo-electronic element can be applied in many different fields; for example, used as polymer LEDs, polymer semiconductor laser, solar cell elements, etc, but not limited to thereto. Please refer to FIG. 5 . FIG. 5 schematically illustrates the cross-sectional view of the light-emitting element of the embodiment of this invention. As illustrated by the figure, according to the light emitting element 2 of this invention includes a substrate 20 , a positive electrode 22 , a hole transporting layer 24 , a light emitting layer 26 and a negative electrode 28 . In practice, the substrate 20 can be made of a transparent glass or other appropriate materials. The positive electrode 22 can be made of a conductive material, such as Indium oxide (ITO). The hole transporting layer 24 can be made of conductive polymeric material (3,4-polyethylenedioxythiophene/polystyrenesulfonate blend, PEDOT/PSS). The light emitting layer 26 is made of polymer molecular film of this invention. Moreover, the negative electrodes 28 can be made of aluminum or other metals. In practice, according to the photo-electronic element of this invention apart from the aforementioned substrates and polymer molecular films, it also includes a protection layer formed on the polymer molecular film to prevent the polymer molecular film from oxidation or wear. Moreover, the photo-electronic element of this invention can include positive and negative electrodes respectively disposed on the substrate and the protection layer. Of course, in practice, the photo-electronic element of this invention can also include other functional layers depending on the situation without being limited by any specific types. Furthermore, apart from the aforementioned designs that utilize the positive and negative electrodes positioned across the thickness of the molecular layer or layers, a parallel version of the electrodes with the positive and negative electrodes situated on a plane parallel to the plane of the molecular film can also be opted. To conclude, the polymer molecular films and the photo-electronic element of this invention increase the light-emitting/power-generating efficiency by the indentation areas formed by stretched conjugated polymers. Moreover, this process can be achieved without taking much time and costing much money, which is highly industrially applicable. Although this invention has been disclosed better as above by the embodiments, they are not intended to limit the scope of this invention. An ordinary skilled person in the art can make any modification and improvements without departing from the spirit and scope of this invention. Therefore, the protection scope of this invention is defined by the appended claims.

The present invention provides a polymer molecular film, a photo-electronic element comprising the same and method for manufacturing the same. The polymer molecular film of the invention is formed via a deformation process on a substrate. Particularly, the polymer molecular film includes a plurality of conjugated polymers, therein at least one of the conjugated polymers has a stretched molecular structure. As a result, the photo-electronic element having said polymer molecular film performs with good lighting or power generating efficiency.

CROSS REFERENCE TO RELATED APPLICATIONS The present application is a continuation of International Patent Application No. PCT/DE01/01581, filed Apr. 25, 2001, designating the United States of America, the entire disclosure of which is incorporated herein by reference. Priority is claimed based on German Patent Application No. 100 20 238.1, filed Apr. 25, 2000; German Patent Application No. 100 20 244.6, filed Apr. 25, 2000; German Patent Application No. 100 21 047.3, filed Apr. 28, 2000; and International Application No. PCT/DE01/00188, filed Jan. 17, 2001. FIELD OF THE INVENTION The invention relates to a progressive spectacle lens as set out in the preamble of claim 1 , the lens producing only slight dynamic distortion. Progressive spectacle lenses (also called varifocal lenses, multifocal lenses etc.) are usually understood to be spectacle lenses having a different (lower) power in the region through which a spectacles wearer views an object located at a great distance—hereunder referred to as a distance portion—than in the region (near portion) through which the spectacles wearer views a near object. Located between the distance portion and the near portion is the so-called progressive zone in which the power of the spectacle lens continuously increases from that of the distance portion to that of the near portion. The magnitude of the power increase is also designated as addition power. As a rule, the distance portion is located in the upper part of the spectacle lens and is designed for viewing “to infinity”, whilst the near portion is located in the lower region and is particularly designed for reading. In spectacles for special applications—those for pilots or for monitor work stations are mentioned as examples—the distance and near portions may also be arranged differently and/or designed for other distances. Furthermore, it is possible for a plurality of near portions and/or distance portions and suitable progressive zones to be present. With progressive spectacle lenses having a constant refractive index it is necessary, in order that the power may increase between the distance portion and the near portion, that the curvature of one or both surfaces continuously change from the distance portion to the near portion. The surfaces of spectacle lenses are usually characterized by the so-called principal radii of curvature R1 and R2 at every point on the surface. (Sometimes also the so-called principal curvatures K1=1/R1 and K2=1/R2 are given instead of the principal radii of curvature.) Together with the refractive index n of the glass material, the principal radii of curvature govern the parameters frequently used for an ophthalmologic characterization of a surface: Surface power=0.5·( n− 1)·(1/ R 1+1/ R 2) Surface astigmatism=( n− 1)·(1/ R 1−1/ R 2). Surface power is the parameter via which an increase of power from the distance portion to the near portion is achieved. Surface astigmatism (more clearly termed cylinder power) is a “troublesome property”, because an astigmatism—inasmuch as an eye does not have an innate astigmatism to be corrected—which exceeds a value of about 0.5 dpt results in an indistinctly perceived image on the retina. BACKGROUND OF THE INVENTION Although any change of the curvature of the surface which is needed to achieve a surface power increase without vision being “disturbed” by surface astigmatism can be attained relatively simply along a (plane or winding) line, considerable “intersections” of surfaces will result alongside this line, leading to a large surface astigmatism which more or less impairs the lens in regions alongside the mentioned line. Furthermore, the strong variation of the prismatic powers results in dynamic distortion effects, the cause of which will be explained in the following: For this, the image of a group G of n pairs of different object points P(x, y) which are located in a plane at a distance s in front of the spectacle lens will be considered. Without limiting the generality, these shall be disposed in the form of a grid having equal spacings. If these object points P(x, y) are imaged by a spectacle lens in such manner that the principal rays pass through a point Z located on the eye side (for example, the center of rotation of the eye or the entrance pupil of the eye), and if the eye-side principal rays intersect a second plane located at a distance r from the spectacle lens, in the following referred to as a projection plane, then a second group B of points results which are the image points of the object points P(x, y) in the projection plane. Generally stated, the spectacle lens performs an imaging A from the object plane onto the projection plane in such manner that A:GB As a rule, the imaged grid is no longer equally spaced, but warped because of a known image defect which is termed “distortion”. When a back-side aperture is present in a path of rays, as is the case with a spectacles wearer, the distortion is cushion-shaped with a single vision positive lens and barrel-shaped with a single vision negative lens. With progressive lenses mixed forms may arise. According to the invention it has been realized that additional effects arise when a temporal change of the distorted grid {right arrow over (ν)} B in the projection plane with translational movements of the object grid {right arrow over (ν)} G are considered. The movement of the grids are represented in a usual manner by vector fields for the velocity {right arrow over (ν)}(x,y). The grid point having the subscript i has the velocity {right arrow over (ν)}(x i ,y i ). It must be remarked that within a particular simply configured region each point P(x, y) in the object plane is imaged onto the projection plane. The limitation to a countable finite number of discrete grid points Pi is made only for the sake of graphical clarity. Thus, n may be finite or infinitely large. When a human eye is confronted with an internally rigid, but moving object grid which has been imaged by a progressive lens, then during the movement the distorted grid will appear to move in the projection plane not only rigidly as a whole. Rather than this, in addition to the expected directed movement caused by the translational movement of the object grid, an undirected component will be observed {right arrow over (ν)} B ={right arrow over (ν)} BDirected +{right arrow over (ν)} BUndirected . In general, the vector field {right arrow over (ν)} B is free neither from divergence nor rotation: There will be regions in which the density of the grid points increases during the movement ∇·{right arrow over (ν)}>0, and others in which it decreases ∇·{right arrow over (ν)}<0: the grid will appear to be subjected to a kneading operation. For rotating movements of glance the following also applies: ∇×{right arrow over (ν)}≠0. The dynamic effects are particularly striking when the directed component is subtracted from the total velocity field. All statements made in the present application concerning track curves or trajectories relate to so-called trajectories in a coordinate system of a particular reference point which is the stationary point. It is expedient to select this marked point to be close to the center of the lens; within the scope of this application and without prejudice to generality, it has the coordinates (0, 0) in the used Cartesian coordinate system which has its origin at the object-side vertex of the spectacle lens. The z axis is directed along the direction of light. The surface-perpendicular vectors of the mentioned planes have only one z component. To describe the movement of the object points arranged in the form of an equally spaced grid, in the following the trajectories of the undirected components of the grid points during a horizontal periodic movement of the object grid will be considered. The track curve of an arbitrarily chosen object point will then be a horizontal straight line, and that of the conjugated image point will be a curve lying in a plane. This curve is characterized by the variation of the prismatic power (or the prismatic secondary effect) along the principal rays generating the curve. The horizontal movements of glance discussed here frequently occur, for example during reading, when driving a car, or when working with a computer. For a single-vision lens having vanishing power the undirected component is equal to zero. The trajectories degenerate to points. The condition of the grid is stationary. In the general case of non-vanishing power, single-vision lenses have the characteristics shown in FIG. 1 . In the left-hand partial illustration the trajectories of a horizontal movement are illustrated for a negative lens having a spherical power of −2.5 dpt, and in the right-hand partial illustration for a positive lens having a spherical power of +2.5 dpt. The stationary point at the center, to which the undirected velocity component of the imaged grid is referred, can be seen. On the right and left of this stationary point are purely horizontal tracks having a vanishing vertical component: with a purely horizontal movement of glance only the horizontal component of the prismatic secondary effects varies along the trajectories in this region; the vertical component—and with it the slope of the track curves—is equal to zero. The length of the plotted trajectories monotonously increases from the stationary point outwards in the radial direction—as also does the prismatic secondary effect. When moving upwards or downwards from the stationary point, an increasing curvature of the track curves may also be discerned. This results because the vertical deflection of the principal rays during a movement varies more strongly than it does further inwards. The position of the centers of curvature correlates with the distortion. An inwardly directed opening of the curves (towards the center of the lens) signifies a barrel-shaped distortion, whilst an outwardly directed opening signifies a cushion-shaped distortion. With progressive spectacle lenses the features of the trajectories for single vision lenses are noticeably changed by the power increase of progressive spectacle lenses. FIGS. 2, 3 , 4 and 5 show trajectories of spectacle lenses on the market; these appear to have been constructed in accordance with the following patent publications: FIG. 2 DE-C-28 14 936 FIG. 3 WO 95/27229 FIG. 4 DE-C-43 42 234 FIG. 5 U.S. Pat. No. 4,606,622 or DE 196 12 177. FIGS. 2 to 5 show right-hand side lenses which are on the market; these have the prescription sph +0.5 dpt, cyl 0, Add 2.0 dpt, Pr 0 (all plots were computed for r=0 mm and s=40 mm). The trajectories of a relatively old progressive spectacle lens shown in FIG. 2 differ from those of all other shown lenses by being, in the bottom half, very short, extremely curved trajectories. A large number are even retrograde, i.e. with the horizontal movements of glance described here, many points appear to move at first with, and later counter to the direction of movement of the objects. This produces serious “swaying sensations”, and the objects appear to be strongly dynamically distorted. For computing a progressive surface in the wearing position, a wearing situation is established. This relates either to a particular user for whom the individual parameters have been specially determined in the respective wearing situation and the progressive surface has been separately computed and fabricated, or to average values as described in DIN 58 208, Part 2. DESCRIPTION OF THE INVENTION The invention is based on the object of further developing a progressive spectacle lens as set out in the preamble of patent claim 1 in such manner that the dynamic distortion which of necessity arises with progressive spectacle lenses has been minimized to the extent that it is no longer felt to be disturbing by a spectacles wearer. An achievement of this object in accordance with the invention is set out in the patent claims. According to the invention, to minimize a dynamic distortion, track curves (trajectories relative to a stationary point at (0, 0)) formed by connecting points of intersection of image-side principal rays passing through a center of rotation of an eye with a projection plane at a distance s from an object-side vertex of the spectacle lens, when horizontally moving objects having coordinates (x−dx, y, s) at a beginning of a movement and (x+dx, y, s) at an end of a movement are imaged through the progressive spectacle lens with r=0 mm, s=−40 mm and dx=35 mm, satisfy the following condition: The absolute value of a difference between a minimum and a maximum y coordinate of a trajectory is smaller than a value H given in the following table: Sph = −4.0 dpt, Add = 1.0 dpt Y = −5 mm X [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.1856 0.1233 0.0917 0.1602 Y = −10 mm X [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.5294 0.3639 0.2681 0.4506 Sph = −4.0 dpt, Add = 2.0 dpt Y = −5 mm X [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.0578 0.0320 0.0032 0.0354 Y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.3601 0.2496 0.1334 0.2765 Sph = −4.0 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.0764 0.0648 0.0919 0.0754 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.1425 0.1088 0.0297 0.0702 Sph = 0.5 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.2399 0.1645 0.1494 0.2164 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.4250 0.2887 0.2684 0.3859 Sph = 0.5 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.4535 0.3059 0.2945 0.4156 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.7102 0.4717 0.4802 0.6734 Sph = 0.5 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.7038 0.4646 0.4508 0.6353 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 1.0653 0.6846 0.7621 1.0458 Sph = +5.0 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 0.8607 0.5828 0.4944 0.7478 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 11.6223 1.2440 1.0375 1.5808 Sph = +5.0 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 1.1921 0.7903 0.7106 1.0740 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 11.3204 1.5119 1.3512 2.0477 Sph = +5.0 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 1.6290 1.0524 0.9997 1.4771 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 H [mm] 10.9637 1.8600 1.8110 2.6494 Alternatively or additionally the arc length of the trajectory may be shorter than the value L given in the following table: Sph = −4.0 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 1.5909 0.9767 0.9266 1.4663 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 1.8098 1.1431 1.1008 1.6803 Sph = −4.0 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 1.6783 1.0503 1.0222 1.5680 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 1.8182 1.1592 1.1855 1.7506 Sph = −4.0 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 1.7463 1.1216 1.1184 1.6751 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 1.8358 1.2066 1.3027 1.8575 Sph = 0.5 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 0.6897 0.4330 0.3721 0.5844 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 0.8672 0.6191 0.4724 0.6851 Sph = 0.5 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 0.9082 0.5691 0.4604 0.7418 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 1.3123 0.9671 0.7530 0.9919 Sph = 0.5 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 1.2493 0.8136 0.6038 0.9704 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 1.9289 1.4589 1.1334 1.4221 Sph = +5.0 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 3.7402 2.2660 2.0300 3.2926 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 17.0238 2.8339 2.3619 3.7662 Sph = +5.0 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 4.1363 2.4717 2.1643 3.6054 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 17.7989 3.2727 2.6076 4.2396 Sph = +5.0 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 4.6985 2.8115 2.3838 3.9886 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 L [mm] 19.1536 3.9888 3.0181 4.7864 Preferably or alternatively the mean gradient of the trajectory may be smaller than the value m given in the following table: Sph = −4.0 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.1504 0.2091 −0.2338 −0.1701 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.3796 0.5037 −0.4935 −0.4012 Sph = −4.0 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.0486 0.0601 0.0086 −0.0375 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.2906 0.4609 −0.3178 −0.2642 Sph = −4.0 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m −0.0698 −0.1448 0.2513 0.0771 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.1408 0.3482 0.0943 −0.0701 Sph = 0.5 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.4286 0.5223 −0.8628 −0.6280 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.5658 0.5608 −0.7568 −0.7614 Sph = 0.5 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.6271 0.7017 −1.2670 −1.0635 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.6496 0.5786 −0.8596 −0.9913 Sph = 0.5 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.7248 0.7639 −1.5740 −1.4026 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.6752 0.5619 −0.9717 −1.1853 Sph = +5.0 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.2821 0.3767 −0.5349 −0.3554 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m −2.0435 0.5602 −0.6714 −0.5836 Sph = +5.0 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.3463 0.4457 −0.6708 −0.4628 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m −1.8039 0.5635 −0.7148 −0.6528 Sph = +5.0 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m 0.4117 0.5009 −0.8530 −0.5979 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 m −1.5811 0.5560 −0.8075 −0.7671 Preferably or alternatively the maximum gradient of the trajectory may be smaller than the value M given in the following table: Sph = −4.0 dpt, Add = 1.0 dpt y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 M 6.7715 0.7457 −0.4274 −0.5090 Sph = −4.0 dpt, Add = 2.0 dpt y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 M 0.2870 0.4403 0.1951 −0.3184 Sph = −4.0 dpt, Add = 3.0 dpt y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 M 0.2179 −0.0023 1.4361 13.7067 Sph = 0.5 dpt, Add = 1.0 dpt y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 M 0.8577 0.7249 −1.1962 −11.5313 Sph = 0.5 dpt, Add = 2.0 dpt y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 M 0.9196 0.9184 −1.3399 −2.8994 Sph = 0.5 dpt, Add = 3.0 dpt y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 M 1.0969 1.0391 −1.9120 −4.9332 Sph = +5.0 dpt, Add = 1.0 dpt y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 M 1.9493 6.1058 −1.9300 −1.7983 Sph = +5.0 dpt, Add = 2.0 dpt y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 M 2.8737 0.8570 −2.1444 −9.0899 Sph = +5.0 dpt, Add = 3.0 dpt y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 M 0.8090 0.7053 −7.1729 −7.8544 Furthermore it is of advantage when preferably or alternatively the x coordinate of the center of the trajectory (half the sum of the minimum and the maximum x coordinate) is smaller than the value xz according to the following table: Sph = −4.0 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −11.3070 −7.5948 7.6229 11.3708 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −11.2191 −7.5330 7.5648 11.2896 Sph = −4.0 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −11.3556 −7.6222 7.6442 11.4147 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −11.3064 −7.5867 7.6028 11.3640 Sph = −4.0 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −11.3994 −7.6451 7.6506 11.4351 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −11.3968 −7.6411 7.6255 11.4162 Sph = 0.5 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −13.7030 −9.1089 9.0833 13.6540 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −13.7736 −9.1541 9.1126 13.7117 Sph = 0.5 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −13.7989 −9.1581 9.1193 13.7287 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −13.9392 −9.2439 9.1778 13.8389 Sph = 0.5 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −14.0239 −9.2904 9.2212 13.9123 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −14.2516 −9.4367 9.3185 14.0854 Sph = +5.0 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −17.3295 −11.3902 11.2874 17.1168 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −12.9054 −11.5804 11.4231 17.3537 Sph = +5.0 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −17.7640 −11.6405 11.5155 17.5159 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −12.8237 −11.8832 11.7004 17.8471 Sph = +5.0 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −18.2779 −11.9499 11.7829 17.9519 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 xz [mm] −12.7057 −12.2789 12.0264 18.3631 Furthermore it is preferred when alternatively or additionally the yx coordinate of the center of the trajectory (half the sum of the minimum and the maximum y coordinate) is smaller than the value yz according to the following table: Sph = −4.0 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −3.9197 −3.9560 −3.9576 −3.9235 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −9.7598 −9.8352 −9.8493 −9.7834 Sph = −4.0 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −3.9663 −3.9993 −3.9887 −3.9561 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −9.9049 −9.9681 −9.9741 −9.9176 Sph = −4.0 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −4.0196 −4.0496 −4.0167 −3.9821 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −10.0654 −10.1113 −10.0993 −10.0540 Sph = 0.5 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −4.5384 −4.5218 −4.5064 −4.5176 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −11.3786 −11.3388 −11.3147 −11.3468 Sph = 0.5 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −4.6100 −4.5783 −4.5450 −4.5643 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −11.5798 −11.5070 −11.4737 −11.5325 Sph = 0.5 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −4.7552 −4.7015 −4.6287 −4.6610 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −11.9218 −11.8067 −11.7550 −11.8468 Sph = +5.0 dpt, Add = 1.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −5.5028 −5.4121 −5.3601 −5.4312 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −18.8206 −13.7124 −13.6144 −13.7883 Sph = +5.0 dpt, Add = 2.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −5.7060 −5.5826 −5.5047 −5.6056 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −18.9496 −14.1505 −14.0444 −14.2803 Sph = +5.0 dpt, Add = 3.0 dpt y = −5 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −5.9781 −5.8024 −5.6754 −5.8106 y = −10 mm x [mm] −30.0000 −20.0000 20.0000 30.0000 yz [mm] −19.0922 −14.6873 −14.5602 −14.8518 In any case it is preferred that inter- or extrapolated values apply to not listed prescriptions. BRIEF DESCRIPTION OF THE DRAWINGS In the following the invention will be described by way of example, without limitation of the general inventive concept, with the aid of embodiments with reference to the drawings to which attention is expressly drawn concerning the disclosure of all details of the invention not described more explicitly in the text. Shown by FIG. 1 are the trajectories for spherical spectacle lenses; FIGS. 2-5 are the trajectories for known progressive spectacle lenses; FIG. 6 are the trajectories for a progressive spectacle lens according to the invention; FIG. 7 is an illustration for comparison of the spectacle lens according to the invention with prior art; and FIG. 8 a are the vertex heights of the progressive surface of a concrete embodiment; FIG. 8 b are the iso-lines of the astigmatic deviation; FIG. 8 c are the iso-lines of the mean “as worn” power; FIG. 8 d are the iso-lines of the surface astigmatism; and FIG. 8 e are the iso-lines of the mean surface power for this embodiment. DESCRIPTION OF EXAMPLES OF EMBODIMENT FIG. 6 shows the trajectories for a spectacle lens according to the invention, having the same optical parameters as the conventional spectacle lenses illustrated in FIGS. 2 to 5 . This illustration already shows that the trajectories are substantially shorter than in prior art and moreover extend considerably more distinctly. This will be set out in greater detail with reference to FIG. 7 . In this illustration the curves for the known progressive spectacle lenses shown in FIGS. 2 to 5 and for the spectacle lens of the invention have been plotted on an enlarged scale, so that the differences between the lenses become much clearer than in the grid plots which serve to provide an overall view and demonstrate the local dependencies in the spectacle lens. FIG. 7 shows that distinct differences of properties such as length of the (individual) track curves, maximum and minimum gradient of the curves, position, etc. exist. These differences are given in the claims in the form of tables for various powers and addition powers. By way of example this comparison will be illustrated by means of the trajectories which can be found by starting from the stationary point and counting (5 downwards, 6 to the left). The object point imaged on the center of the trajectory thus has the coordinates (−30 mm, −25 mm, −40 mm). It is moved horizontally by a total of 70 mm, i.e. from the coordinate x=−65 mm to x=5 mm. The spectacle lens illustrated in FIG. 2 is the “curl” in the top right-hand corner of FIG. 7 . For the spectacle lens illustrated in FIG. 3 the trajectory is distinctly longer than for the spectacle lens illustrated in FIG. 2, however it is elongate and thus does not show the disturbing effects of being retrograde. The curve rises monotonously everywhere. When the gaze is kept fixed on a moving object point whilst the head is kept at rest (or when a corresponding head movement is made whilst the object is stationary), the spectacles wearer must monotonously raise or lower his glance in accordance with the curve. This uniform monotonous movement of the glance is experienced as being more agreeable than that with the retrograde strong curve of the spectacle lens illustrated in FIG. 2 . The spectacle lens illustrated in FIG. 4 has agreeable properties similar to those of the spectacle lens according to FIG. 3 . However, it must be noted that the maximum gradient is larger than that according to FIG. 3 . To a wearer making a movement of glance, a peripheral object will appear to fall or rise more rapidly in the same time than to a wearer of a spectacle lens according to FIG. 3 . If adjacent trajectories are examined, it will become clear that objects also become distorted much more rapidly and that the distortion therefore becomes more noticeable. Apart from the gradient being greater, it will be noticed that the curve is located at a distinctly lower position than the others. The necessary lowering of the glance is greater than with the other products. The spectacle lens according to FIG. 5 also has a relatively large gradient. The spectacle lens of the invention as illustrated in FIG. 6 avoids all the described disadvantages of the other lenses during horizontal movements of glance: The relative trajectories are here distinguished by short length, monotony, and smallest mean and smallest maximum gradient simultaneously with an only slight lowering of the glance. The progressive spectacle lens of the invention is thus distinguished by the dynamic distortion with horizontal movements of glance and the swaying movements correlated therewith are as small as possible. The concrete example of embodiment illustrated in FIG. 8 has a spherical power (mean “as worn” power) of −1 dpt and an addition power A of 2 dpt at the distance reference point. An astigmatic prescription is not present. In all Figures the abscissa (x axis) is the horizontal axis and the ordinate (y axis) is the vertical axis in the wearing position. The distance and the near reference point are each represented by a circle in FIGS. 8 b-e , the centration point being designated by a cross.—their positions may be seen in the Figures. Furthermore, the course of the principal line has been plotted. The partial FIG. 8 a indicates the vertex heights of the progressive eye-side surface for the embodiment. Vertex height is understood to be the distance of a point having the coordinates x and y (horizontal and vertical axis, respectively, in the wearing position of the spectacle lens) from the tangential plane of the surface vertex. In the Table each left-hand column shows the y values (from −20 to +20 mm) and the top line from column 2 onwards shows the x values (from −20 to +20 mm). The vertex heights are also given in millimeters. The value 0 means that no vertex height is given for these x,y coordinates. The partial FIG. 8 b shows the astigmatic deviation within a circle having a radius of 30 mm around a point lying 4 mm below the so-called centration point. The astigmatic deviation is the “residual astigmatism” of the system spectacle lens/eye and is shown by so-called iso-lines beginning with the iso-line 0.25 dpt. The iso-lines indicate the deviation of the astigmatism with respect to magnitude and cylinder axis from the cylindrical prescription—which in the case of an astigmatism-free eye is 0 dpt. The partial FIG. 8 c shows the corresponding iso-lines for the mean “as worn” power of this embodiment. The mean “as worn” power D is the mean value of the reciprocals of the image side focal intercepts S′1 and S′2 minus the object distance which is the object side focal intercept S, as given by D= 0.5·( S′ 1 +S′ 2)− S and is also illustrated in the form of so-called iso-lines beginning with the iso-line 0.75 dpt. In corresponding manner the iso-lines of the surface data, namely the surface astigmatism and the mean surface power, are shown in partial FIGS. 8 d and 8 e . For a definition of these surface data attention is drawn to the introductory explanations. The embodiment shown in FIG. 8 has the following individualized conditions of wearing: D1x 4.55 D1y 4.55 n 1.597 d 1.59 DRP 1.0 PD 63 HSA 15 Pantoscopic Angle 0 wherein: D1x is the surface power of the front surface in x direction (dpt) D1y is the surface power of the front surface in y direction (dpt) n is the refractive index of the glass material d is the center thickness of the lens in mm DRP is the prism thinning in cm/m PD is the interpupillary distance in mm HSA is the vertex distance in mm. The pantoscopic angle of the spectacle lens is given in degrees. Of course, the features of the invention may also be used for computing and manufacturing spectacle lenses having two progressive surfaces and/or having (in addition) a varying refractive index.

Described is a spectacle lens comprising a region (distance portion) designed for viewing at large distances and in particular “to infinity”; a region (near portion) designed for viewing at short distances and in particular “reading distances”; and a progressive zone disposed between the distance portion and the near portion, in which the power of the spectacle lens increases from a value at a distance reference point located in the distance portion to a value at a near reference point located in the near portion along a curve (principal line) veering towards the nose. The invention is rendered distinct by a feature, amongst others, according to which trajectories of motion fulfill specific conditions.

FIELD OF THE INVENTION [0001] The present invention relates to the field of plastic fabrication and uses. More specifically, the present invention relates to multilayered plastic films that have antifogging and antimicrobial properties. BACKGROUND [0002] It is known that many thermoplastic polymer packaging materials, such as films, coatings, sheets, bags, and the like, with suitable strength and flexibility are used to enclose perishable foods, fruits, raw meats, daily dishes and vegetables. These packaging materials tend to discolor and fog during extended storage. Because of this, polymer packaging materials have to possess the following characteristics: (1) suitable thickness and cohesive properties for packaging, (2) high antifogging properties, i.e. the films do not accumulate water droplets on the surface of the material, (3) high mechanical strength at break, (4) appropriate slip properties, (5) excellent optical characteristics, such as gloss and transparency, and (6) sealability under heat. [0003] There is a high demand in the packaging food industry, agriculture, industrial markets, flower wrapping trade, and the like for biaxially oriented thin antifogging films of different types that can be used for both food-wrapping and agricultural applications. The antifogging and antimicrobial films reduce the growth of living contaminants (such as bacteria and molds) and ensure that any condensation of water vapor occurs as an uniform, invisible, layer of water rather than as a series of individual droplets which are not only aesthetically undesirable but produce damaging effects. [0004] The several goals of these films are: (1) to ensure that the polymer thin films retain their transparency so that the packaged contents are clearly visible and so that there is maximum light transmission into the enclosure, (2) to protect the packaged food products from undesired degradation that may be caused by the droplets of water, (3) to prevent large drops of condensed water from falling onto young plants, increasing the possibility of damage and disease, (4) to prevent plant “burning” caused by large drops of water lensing (concentrating and focusing) solar radiation onto the contents of the package, (5) to provide antimicrobial properties and (6) to provide prolonged shelf life by preventing the growth of the certain bacteria. [0005] Currently, antifogging films (also known as antifog or antimist films) are produced by adding or coating various types of organic antifogging additives, such as ethoxylated sorbitan ester, glyceride fatty acid ester, glycerol stearate (or monostearate), glycerol oleate and sorbitan ester, and the like, to conventional film forming polymers, such as polyolefins, flexible vinyl chloride polymers, oriented styrene polymers, polyesters, ethylene-vinyl acetate copolymers, and the like. [0006] There are a number of available patent publications related to antifogging polymer films obtained by using different types of thermoplastic film-forming/antifogging additives as discussed below. These patents relate to systems such as a biaxially stretched film with a base of an olefin polymer resin composition containing ethylene-propylene copolymer and 0.5% of polyethylene glycol stearyl ether, olefin polymer/fatty acid monoester of polyhydric alcohol (or alkaline metal salt of a diester of sulfosuccinic acid), polyolefin/ethylene oxide (or monoglyceride of a fatty acid), polystyrene/alkyl phenyl polyethylene glycol ether (of fatty alcohol sulfate) base coating, polyethylene/polyhydric alcohol esters or metal salts of either saturated or unsaturated monocarboxylic fatty acids, ethylene polymer and polybutene blend/glyceride with acyl group, and ethylene-acrylic acid (or ethyl acrylate and/or vinyl acetate) copolymers or low density ethylene polymers/alkyl phenyl polyethylene glycol ethers or alkoxylated alkyl phenol. But all of these patents suffer from one or more of the following disadvantages such as higher haze values, low values of sheen, higher transverse or longitudinal shrinkage, and poor antifogging properties. [0007] More specifically, in U.S. Pat. No. 4,066,811, there is disclosed raw tubular polyolefin films with suitable orientation determined by heat shrinkage, containing ethylene-vinyl acetate copolymer, polyethylene, polypropylene or mixtures thereof, polyalkylene ether polyol and non-ionic surfactant-polyhydric alcohol ester derivatives of fatty acids. In the above patent, the determination of antifogging properties of the subject film was according to the following measurements: (1) no water droplets were present on the surface and water was in a uniform layer, (2) large water droplets locally were adhered or there was unevenness in the state of any adhering water droplets, and (3) fine water droplets adhered to the whole surface. [0008] Other recently published patented inventions, such as JP Pat. 09-104,092, relate to various polymer compositions, sheets, and films having fog resistant properties. Disclosed therein are antifogging sheets comprising weather-resistant polycarbonate based films, hot-melt poly(methylmethacrylate) films containing a benzotriazol UV-absorber, and cellulose films containing a diethyl phthalate plasticizer, to form a flat or wavy laminated panel allegedly providing good weather and moisture resistant adhesion. [0009] Antifogging polypropylene lids with smooth handling properties, such as disclosed in JP Pat. 09-76,339, were prepared by thermal formation of polypropylene sheets, where the interior faces of the lids exhibit antifogging property and the exterior faces have a friction coefficient of 0.01-0.7. The plastic of these lids was stretched in the machine direction, coated on the exterior face with poly(dimethylsiloxane) and on the interior face with sugar fatty ester emulsion, and thermal formed into a lid showing no noise when removed from their stack. [0010] Plastic sheets having anisotropic surface characteristics, including fogging and adhesion properties, are disclosed in JP Pat. 09-85,847 and comprise alternating strips of nylon 6-12 and ethylene-methacrylic acid copolymer. [0011] There are antifogging laminated films for agricultural uses that use a polyolefin resin middle layer. This layer frequently consists of high density polyethylene and synthetic rubber with external layers consisting of antifogging agents. One laminate, disclosed in JP Pat. 0994,930, comprises an ethylene-vinyl acetate copolymer middle layer, uses KFG 561 as an antifogging agent, and showed good blocking resistance, mechanical strength and fogging presentation (45° C. water for 45 days or 0° C. environment and 20° C. water for 24 hours). [0012] Other agricultural antifogging films, such as the ones disclosed in JP Pat. 09-95,545, were prepared using olefin copolymer compositions containing sulfonated olefin copolymers, ethylene-C 3-12 olefins, and ethylene-acrylic copolymers. The olefin copolymers were synthesized by polymerization of olefins in the presence of metallocene (Zr) catalyst containing silica and methylaluminoxane. More specifically, a transparent antifogging film was prepared from a mixture of 80% sulfonated olefin polymer (reaction product of butane sulfonate with ethylene-acrylic copolymer) and 20% of ethylene-hexene-1 copolymer which was polymerized in the presence of a catalyst system containing silica, methylaluminoxane, bis(1,3-n-butylmethyl cyclopentadienyl)zirconium dichloride and triisobutylaluminum. [0013] JP Pat. 09-77,938 discloses a polymer composition with good sliding properties that comprises 10-60% of graft copolymers manufactured by grafting an elastomer with ≧1 layers of antifogging agent KFG 561. The resultant laminate used ethylene-vinyl alcohol copolymer as a middle layer, 20% of hydrogenated butadiene-styrene elastomer as an inner layer, and 10% of the said elastomer outer layer comprises a fire retardant agent and showed good blocking resistance, mechanical properties, dust, and fogging presentation. [0014] Fluoropolymer films with wetting ≧35 dyn cm, as disclosed in JP Pat. 09-136,980, were mixed with antifogging agents comprising water-thinned acrylic polymer emulsions, such as ethyl acrylate-2-hydroxyethylmethacrylate-2-hydroxymethacryloxybenzophonone-methyl methacrylate copolymer, and colloidal SiO 2 . Films prepared according to this patent showed reasonably good antifogging property for 7 months. [0015] Two Japanese patent inventions, JP Pat. 09-165,178 and JP Pat. 09-165,447, disclose heat-aging and light-resistant propylene polymer compositions causing no fogging of glass for use in automotive interiors. These compositions contain (A) crystalline polypropylene, (B) inorganic filler, such as TiO 2 , (C) ethylene-propylene elastomer, and (D) conventional stabilizers, antioxidants, antiblocking agents, and other additives such as epoxy resins, hydroxyl-containing low molecular weight polyolefins, polyethylene waxes, and anionic surfactants. Plates prepared from this composition by kneading, pelleting, and injection molding show 150° C. oven life for 320 hours. The plate and glass plate were left in a sealed container at 120° C. for 20 hours and showed a haze of the glass of 0.8%. [0016] In another patent entitled “Fog-Resistant Heat-sealable Film”, U.S. Pat. No. 4,341,825, there is disclosed a transparent, heat-sealable, laminated film that has a first layer of a difficulty heat-sealable polymer, such as an axially oriented polyethylene terephthalate film with 0.002-0.006 cm thickness, and a second layer of a readily heat-sealable polymer, such as low density polyethylene and copolymers of ethylene with acrylic acid, ethyl acrylate and vinyl acetate, chemically interfacially joined to the first film layer. The said second film layer comprises 0.2-0.7% of an alkyl phenyl polyethylene glycol ether of the formula, R—C 6 H 4 —O—(CH 2 ) n OR′—OH, where R-alkyl C 10-15 and alkylene C 4-10 as an antifogging agent. The resulting laminated film is then heated to 130° C. and exposed to UV-light through the second film layer for a time and at an intensity sufficient to cause the formulation of the chemically interfacial bond between the two layers. The film obtained resists the formation of fog when utilized to package refrigerated foods. However, the disadvantages of this invention can be noted as the following: (1) the subject film comprises two layers containing non-oriented ethylene polymers, (2) the subject film has a high thickness, (3) the subject film has a high content of antifogging agents as compared with more conventional polymer fog-resistant films, and (4) the antifogging agents used in the subject film were synthesized by reaction of alkyl phenol with polyethylene oxides. In this case, the trace of the phenol will be present in the product synthesized. This can limit the use of this specific additive in the food packaging industry. [0017] Another patent, entitled “Fog-Resistant Olefin Polymer Films”, U.S. Pat. No. 4,486,552, discloses a film-forming composition for making packaging films that are resistant to fogging, especially when employed as a packaging film for moist products. The subject film of this patent comprises an ethylene polymer, especially a linear low density polyethylene, and 0.5-2.0% of antifogging agents, such as an ethoxylated alkyl phenol along with a mixed mono-, di-, and/or triglyceride, a polyoxyalkylene fatty acid ester or various combinations of said additives. The mixing of the antifogging agents into the ethylene polymers, which can be LDPE, LLDPE, HDPE, ethylene-octene-1, or blends or alloys of said olefin polymers, is done by mixing the antifogging agents into molten polymer by commonly used techniques, such as roll-milling, mixing in a Banbury type mixer, mixing in an extruder barrel, or the like. The subject film was formulated as 0.015 mm on a cast film unit at 260° C. melt temperature and chill roll temperature of 18° C. It is noted that the films prepared according to this patent have a relatively high fog resistance when compared with commercially available plasticized poly(vinyl chloride) films, such as the one disclosed in U.S. Pat. No. 4,072,790. Further, other high qualities are produced, such as improved transparency (64.3 against 5.0 for PVC), gloss (95.9 against 89.0), haze (1.0% against 2.0%), lower heat seal range (121-127° C. against 149-177° C.), and overall toughness, as compared to PVC films. However, it was shown that the antifogging agents used in this patent exude to the surface of the film within approximately 48 hours after fabrication. The subject films of this patent have the following disadvantages: (1) the films are not multi-layered and biaxially oriented, (2) the films have a high thickness and high density resulting in a low yield, (3) there is a low heat-sealing temperature, (4) there are low values of surface and mechanical characteristics, the film surfaces are not treated by corona discharge, and (5) the film comprises relatively high concentrations of antifogging agents used in the polymer composition. [0018] U.S. Pat. No. 4,876,146 and 4,956,209, disclose “Anti-fogging Multi-layered Film and Bag Produced Therefrom for Packaging Vegetables and Fruits”. These patents describe biaxially oriented and multilayered antifogging polyolefin films useful for packaging fresh vegetables and fruits comprising: (A) a 4-100 μm base layer formed from polypropylene or ethylene-(5%)-propylene popolymer or ethylene-vinyl acetate (acrylic acid or styrene) copolymer; and (B) one or two surface layers that are 0.3-8.0 μm thick and having heat-sealable properties resulting from a (1:1) mixture of propylene-butene-1 (18%) and ethylene (3.5%)-butene-1 copolymers containing 0.3-3.0% antifogging agent such as higher fatty acid ester of monoglyceride (or alkyldialcoholoamide, polyalkylene glycol, polyalkylene glycol alkylphenol ether). There are also other conventional additives, such as antistatic and lubricating agents. In accordance with said patent it is possible to incorporate the antifogging agent only in a base layer of the film so that the antifogging agent migrates to and diffuses into the surface layer(s) after laminating the layers. This migration and incorporation of the antifogging agents into the surface layers provides the antifogging property important to the surface layer. Antifogging properties were observed, the film was formed as a bag and “Shtitake” mushroom were enclosed in the bag; the temperature was varied twice per day with a rise and drop between 20 and 40° C.; the result was observed after 1 day. There was little fogging, discoloration, and the measured surface tensions were 38-42 dyne/cm. The disadvantages of the films prepared in accordance with said patent included: (1) high values of haze (3.1%), (2) low values of sheen (86.6%), (3) coloring agent in the film does not comply with food contact standards of the U.S. FDA, (4) the identification of fogging properties used a non-effective method, (5) the films had low performances as antifogging surfaces, i.e. discontinuous film of water is observed on the surface, (6) E-P-B terpolymer is not used in the surface layers, (7) ethylene-vinyl acetate copolymers are used in the base layer and most probably for improvement of barrier properties of films, and (8) present patent is limited to using 2-3 layered films. U.S. Pat. No. 4,876,146 and 4,956,209 which are hereby incorporated by reference. [0019] All of the patents previously mentioned above, however, suffer from not having antimicrobial properties. [0020] In the recent years, essentially growing trend is the use of various bioactive agents, including predominantly ecologically pure metal-containing biocides in polymer production industries for preparation of antimicrobial, antibacterial and antifungal polymer materials such as films, sheets, coatings, plastics, fibers, composits, etc. The number of patent publications in this field have increased in recent years. The following references have attempted to address antimicrobial films: (1) U.S. Pat. No. 4,938,955, 1990 discloses an antibiotic resin composition comprising at least one antibiotic zeolite of which ion-exchangable ions are partially or completely replaced with ammonium ions (5-15%) and antibiotic metal ions (Ag + of 1-15%), at least one discoloration inhibitor such as benzotriazole, oxalide, anilide, salicylic acid, phosphous, sulfur, etc. compounds and at least one polymer resin (this composition exhibits antibiotic property and does not discolour with time, and can be employed to form a variety of products which require antibacterial and/or antifungus properties); (2) Transparent bactericidal multilayer sheets with haze <5% comprise a crystalline thermoplastic resin containing 0.05-5 phr granular zeolite containing bactericidal metal ions in a sheet comprised polypropylene containing 0.5% Bacterikiller BM 103 (zeolite A containing 3.5% Ag) [JP Pat. 04,275,142 (1992), Chisso Co., Japan]; (3) Antibacterial polyolefin compositions with inhibiting effects on the growth of bacteria and moulds contain polyolefins and 2-pyridinethiol 1-oxide and its metal (Zn) salts or other organic biocides (polypropylene 100, 2-(4-thioazolyl) benzimidazole 0.25 and Zn 2-pyridinethiol 1-oxide 0.25 part were roll kneaded at 230° C. and then hot pressed at 220° C. to give a 2 mm sheet, which completely inhibited of the growth of Aspergillus niger, Penicillium citrinium, Chaetomium globosum, Aurebasidium dulllans, and Gliocladium virens at 28° C. for 28 days) [JP Pat. 04,270,742 (1992), Shinto Paint Co. Ltd., Japan]; (4) Antibacterial heat-resistant polyolefin compositions comprising polyolefins (polypropylene)100, bactericidal metal ions (Ag, Cu, Zn and/or Sn ions supported on zeolites) 0.01-1.5, dimethylsiloxane oil 0.01-0.2, and aluminium borate whisker (9Al 2 O 3 .2B 2 O 3 ) 0.01-0.1 part showed good antibacterial action as tested against colon bacilli [JP Pat. 04,363,346 (1992), Tonen Kakagu Kk., Japan]; (5) JP 04,13,733 (1992) discloses antibacterial films for packaging chemicals and food which were prepared by treating one or two surfaces of films containing aluminosilicic acid salts with electrical corona (a composition containing 2 parts zeolite A (Ag content 6.7%, NH 4 content 0.5%) and polyamide (6-nylon 66 copolymer) were together extruded and exposed to electrical corona for 0.2-10 s to give an antibacterial film with good adhesion to ham, versus poor adhesion for the film not treated with said corona); (6) U.S. Pat. No. 5,614,568, 1995 (Mawatari, M., et al., Japan Synthetic Rubber Co., Ltd., Tokyo) claimed an antibacterial resin comprising (A) 100 parts by weight of aromatic alkenyl resin, specifically styrene resin, (B) 0.01-30 parts of an inorganic metal compound or a porous structure substrate which has been injected to ion-exchange with a metal ion selected from the group consisting Ag, Zn, Hg, Sn, Pb, Cd, Cr, Co, Ni, Mg, Fe, Sb and Ba, and (C) 0.01-30 parts of a polyethylene comprising —COOH, —COOM(salts), —OH, —COOR, and epoxy, anhydride and amine functional groups, a polypropylene comprising said selective functional groups with molecular weight 10000-30000; (7) Japan Chem. Ind. Co. (JP Pat. 09,176,370, 1997) discloses an antimicrobial injection-moldable polypropylene composition showing no discoloration or degradation during processing, storage and uses contain 0.2 phr of liquid paraffin, 1.0 phr of the mixture of inorganic compounds Ag 0.15 Na 0.5 H 0.35 , Zr 2 (PO 4 ) 3 and Mg 0.7 Al 0.3 O 1.15 which was used as an antimicrobial agent; (8) Polyethylene terephthalate films coated with thin Ag, Cu and Ti-layers by sputtering treatments have high antibacterial activity. The reducing in bacteria values of almost 100% were determined by the SEK Shake Flask Method and the Contacted Film Method [S. Kubota, et al, Bakin Bobai, 25 (7), 393 (1997); Chem. Abstr., 127, 122386s (1997)]; (9) Tokuda, et al, [JP Pat. 09,136,973, 1997] describes bactericidal packaging films comprising thermoplastic resins or blends on the base of PE, PP, PVC, polyesters and/or PS and calcined powder ceramics containing 40-60% of SiO 2 , 20-30% of Al 2 O 3 , 4-8% of ZnO, 2-5% TiO 2 and 0.1-1.0% Ag or Cu salts as an antibacterial agent (these films were prepared by mixing above ceramics with said polymers and forming into films or by spreading or printing above ceramic-containing resins on resin base films); (10) JP Pat. 09,123,264 (1997) discloses antibacterial decorative sheets and manufacture of decorative moldings (these sheets were prepared by shaking colored base sheets with thermosetting diallyl phthalate resin composition containing 0.5% of Ag/Zr phthalate, Ag tripolyphosphate, Ag hydroxyapatite, and/or (Ag/Ca) 3 phosphate. A printed paper sheet was hot pressed with a moldable polymer composition to form a waterproof pan for bathroom uses); (11) Bactericide-containing abrasive agents and resin moldings for video and arcade games comprise a thermoplastic resin (98% of polycarbonate) incorporated with fillers (1%) and bactericides (Ag-containing zeolite, Bactekiller) or bactericide-treated powders (1%) [Sumitomo Elect. Ind. Ltd., JP Pat. 09,77,880, 1997]; (12) JP Pat. 09,77,042 (1997) releases to antimicrobial synthetic resin containers for preserving drinking water (this container is prepared using synthetic resins with Ag-containing glass particles that release adequate amount of microbiocidal silver ions (Ag + ) into the water where growth of bacteria or fungi in the drinking water is prevented by these ions); (13) JP Pat. 09,002,517, 1998 [Taisho Pharmaceutical Co. Ltd. (Tokyo, Japan)] discloses a process for making a bottle and cap with antibacterial properties on their inner contact surfaces. Antimicrobial zeolite power (1 to 5% by weight) containing microbiocidal Ag, Zn and Cu ions is mixed with thermoplastic resins such as ethylene-vinylacetate copolymer, polypropylene and polyethylene (the zeolite is dispersed throughout the bottle and is present on both inner and outer surfaces and can also be used for both cap and membrane seal); (14) Polypropylene plastic table wares contain an antimicrobial agent (Amenitor) (JP Pat. 09,108,084, 1997); (15) Bactericide power (Bactekiller) or bactericide-treated power containing adhesive agent and resin moldings for video areade games were described (Chem. Abstr., 127, 35460t, 1997); (16) Silver (Ag)-zeolite antimicrobial agents for protection of the plastic films from various microorganisms were manufactered by Michubusi Co. Bactericide ceramic power containing 0.1-1.0% Ag or Cu, 2-5% TiO 2 , 4-8% ZnO or MnO 2 , 20-30% Al 2 O 3 and 40-60% SiC or SiO 2 was recommended to use in the varoious thermoplastic composition (polyolefin, polystyrene, polyesters, etc.), resins and binders [T. Ishitaki, High Polym. Japan, 39(10), 744 (1990); Y. Kajiura, Jidosha Gijutsu, 51(5), 34 (1997); JP Pat. 09,136,973 (1997)]; (17) Antimicrobial activities of some new coordination polymers were also discribed by Patel, et al. [B. Patel and M. Mohon, J. Polym. Mater., 13(4), 261 (1996)]. [0021] However, all these publications are related to the preparation and use of various antimicrobial polymer materials including non-orientated and non-multilayered polymer films, sheet, etc. containing bioactive metal ions. Thus the above patents describe inventions are essentially different from the present patent invention which is concerned with preparation of semi- and biaxialy oriented and multilayered antimicrobial thin films containing Ag + -containing polymeric bioactive agent only in the skin layer and having high physico-mechanical, thermal and antimicrobial properties. Another distinctive feature of these films is possibility of their use in the food packaging applications, where anti-fogging properties are required. [0022] Several Firms such as Taisho Pharmaceutical Co. Ltd. (Tokyo, Japan), Kanebo Chemical Industries, Ltd. (Osaka, Japan), M. A. Hanna Company (ISA, Neutrabac™ Antibacterial Masterbatch), Wells Plastic Ltd. (Staffordshire, UK), etc. have already started to manufacture organic and inorganic antibacterial agents and various antimicrobial Masterbatches for use in thermoplastic polymer compositions. [0023] Many organic and organoelement compounds having high biological activities are also used in polymer film-forming composition systems [Z. M. Rzaev, CHEMTECH, (1),58 (1976); Z. M. Rzaev et al., England Pat. 1,270,922 (1972); U.S. Pat. No. 4,261,914 (1981); U.S. Pat. No. 4,314,851 (1982); Z. M. Rzaev et al., Bioresistant Organotin Polymers, Chemistry, Moscow, 1996 (Russ.)]. Thus, (1) “ICI Biocides” Firm (UK) prepared and patented new water soluble biocides on the base of isothioazolione useful for the effective preservation of polymer resins, specially aqueous-based paints from bio-destruction with microorganisms in the stage of synthesis, storage and uses of these materials [C. L. P. Eacoff, Orient. J. Oil and Colour Chem. Assoc., 74 (9), 322 (1991)]; (2) Polen Kagaku Sangyo K.K. [JP Pat. 09,169,073, 1997] discloses antibacterial and antifungal sheets laminated with low expanded olefin polymer (such as HDPE) compositions containing 0.1-1.0% 2-(4-thioazolyl)benzimidazole as an antibacterial and antifungal agent showing good deep drawability; (3) Antimicrobial rubber articles contain ammonium salt of chlorohexidine as an antimicrobial agent [UK Pat. 8,919,152 (1990)]; (4) Biocide Cl-containing polyketones having antibacterial activity against selected yeast, fungi,—and bacteria were prepared by Fiedel-Graft [Friedel-Craft] reaction of o-cresol with chloroacetyl chloride, dichloromethane and dichloroethane in the presence of anhydrous AlCl 3 as a catalyst in nitrobenzene as solvent [B. T. Petel, et al, Orient. J. Chem., 13 (1), 83 (1997); Chem. Abstr., 127, 136122q (1997)]; (5) Polyethylene four-layered film was coated with mixture of allyl isothiocyanate (as a biocide), polyfunctional isocyanate, polyols and dibutyltin laurate (as a catalyst) to give a multilayered film with polyethylene outer layer having antibacterial activity [JP Pat. 09,151,317, 1997]; (6) Matsukawa Electric Works, Ltd. (Japan) was disclosed a method of preparing plastic table wares (plastic bowl) containing antimicrobial agents (Amenitop) by moulding. The core potion is formed with a polypropylene resin and this is coated with another polypropylene containing a said antimicrobial agent; (7) Kyowa Co. Ltd. [JP Pat. 09,135,716, 1997] patented the gas-permeable and antimicrobial bags for the medical application. These bags were prepared from cushion bases consisting open-celled polymer foams and bactericide-containing hydrophobic noncircular fiber; (8) p-Hydroxy butylbenzoate [JP Pat. 63,173,723 (1988)], 2-(4′-thiazolyl)-benzimidazole [U.S. Pat. No. 4,008,351 (1977)], Pt-vinylsiloxane complex [JP Pat. 04,202,313 (1993)], polymeric iodine complexes [U.S. Pat. No. 3,907,720 (1975), phosphate esters [U.S. Pat. No. 3,888,978 (1975), U.S. Pat. No. 3,991,187 (1976), U.S. Pat. No. 4,661,477 (1987), U.S. Pat. No. 4,935,232 (1990)] and 2,3,5,6-tetrachloromethylsulfonylpyridine (for preparation antibacterial styrene type resin compositions) [JP Pat. 07,82,440 (1995)] have also been recommended for use as bactericide and antimicrobial agent in the various polymer compositions, film and sheets. [0024] There are a number of patents disclosing various polymer composits, thermoplastic fibers, sheets, coatings, films, etc. having biological activity toward different type of microorganisms [Shima et al., U.S. Pat. No. 4,000,102, 1976; Dell et al., U.S. Pat. No. 4,584,192, 1986; Fink et al., U.S. Pat. No. 4,751,141, 1988; Gillete et al., U.S. Pat. No. 5,152,946, 1992; Grighton et al., U.S. Pat. No. 5,246,659, 1993; 5,104,306, Apr. 14, 1992]. For example, (1) U.S. Pat. No. 5,178,495, 1993 discloses a polymeric film with biocide. A multi-ply film has been developed that includes a biocide in at least one the film layers. Said biocide mixed with the thermoplastic prior to extrusion of the sheet. This sheet with biocide can be used to construct water containment facilities for drinking water, fish farms and industrial use and can be used as a covering for water tanks or equipment in environments that promote microbial growth at the surface of the film; (2) U.S. Pat. No. 5,777,010 1998 (Nohr R. S., et al., Kimberly-Clark Worlwide, Inc., Neenah, Wis.) discloses melt-extrudable composition containing antimicrobial siloxane quaternary ammonium salts. These compositions which includes a thermoplastic polyolefin and a siloxane quaternary ammonium salt additive. Upon melt extruding the thermoplastic composition to form fibers and non-woven webs, or other shaped artides, the surfaces of such shaped articles exhibit antimicrobial properties. (3) Early, antimicrobial siloxane quaternary ammonium salts were patented [U.S. Pat. No. 5,567,372, 1994 and U.S. Pat. No. 5.569,732, 1994] and published [ Nohr R. S., et al., J. Biomed. Sci., Polym. Ed., 5(6), 607 (1994)] (U.S. Pat. No. 5,567,372, 1994 discloses a method of preparing a non-woven web containing antimicrobial siloxane quaternary ammonium salts); (4) U.S. Pat. No. 5,527,570 [Addeo, A., et al.,1996, Centro Sviluppo Settori Impiego SRL, Milan, Italy)] relates to a multilayer and multifunctional packaging elements having high-absorption activity toward aqueous liquid substances as well as barrier properties toward gases such as oxygen and carbon dioxide are prepared by thermoforming (Each layer comprises a polymeric thermoplastic material. Intermediate layer of this packaging element may also contain antibacterial agents); (5) U.S. Pat. No. 5,142,010, 1992 (Olstein, A. D. et al., H. B. Fuller Licensing & Financing Inc., Wilmington, Del.) discloses polymeric biocidal agents containing carboxyl groups, fluorene substitute and alkyl C 1-20 groups, and any bioactive naturally occurring amino-acid chain (the resulting polymers are disclosed to be useful in any variety of applications requiring an antimicrobial agent or an active sanitizer or disinfectant including films, coatings and adhesives, as well as also being useful in medial, food preparation and personal care product applications; (6) Describes an antimicrobial film-forming compositions containing bioactive polymers (homo-, co- and terpolymers of monomers containing pyran groups) having pendant pyran groups [Greenwald R. B. et al., U.S. Pat. No. 5,108,740, 1992, Ecolab Inc., St. Paul, Minn.] (this publication describes a liquid composition that yields an abrasion resistant polymeric film on a surface that provides extended protection from microbial growth through slow release of a potent antimicrobial agent). [0025] As evident from the above described patent publications, there is relatively small number of patent publications describing polyolefin based, in particular, mono- and biaxially oriented polyolefin based films, and all of patent publications suffer from one or more of the following properties: not being multilayered and oriented polyolefin non-opaque films, not being heat-sealable; not having antimicrobial properties using thin films containing Ag + -containing polymeric bioactive agent only in the skin layer, not having antifogging properties. SUMMARY [0026] It is an object of the present invention to design and prepare a multilayer structure (having at least an antifogging and antimicrobial skin layer (A)/a core layer (C)/an outer layer (E) structure) for semi and biaxially oriented polyolefin based antifogging films having advantageous properties as compared with known and commercial films such as low values of haze, high values of sheen, lower longitudinal and transverse shrinkage, which provides high dimensional stability, and excellent antifogging and antimicrobial properties. Preferably the antifogging and antimicrobial skin layer (A) is electrical corona or flame treated. Electrical corona or flame treatment of the the outer layer (E) may enhance ink anchorage and increase the printability of this layer. Preferably, the films comprise an inner (B) layer between the antifogging and antimicrobial skin layer (A) and the core layer (C). More preferably, the inner (B) layer has the same composition as antifogging and antimicrobial skin layer (A) without the antimicrobial additives. Preferably, the films may comprise a second inner (D) layer between the outer (E) layer and the core layer (C). More preferably, the second inner (D) layer has a preferred composition of 100 percent (%) E-P-B terpolymer. [0027] Antimicrobial and antifogging ≧3 layers polymer films with preferable A/C/E structure useful for the food, medicine and agriculture applications as well as for other general packaging and non-traditional special applications. More preferably, antimicrobial and antifogging films having a A/B/C/E structure. Most preferably, biaxially oriented polypropylene films having symetrical structure A/B/C/D/E, where two outer layers A and E are having antimicrobial and antifogging properties and heatsealable and two intermediate layers B and D are made of E-P random copolymers or E-P-B terpolymers, with or without antifogging agents. [0028] A preferred embodiment of antifogging and antimicrobial skin layer (A) comprises the following compositions: polypropylene greater than or equal to 1 percent (wt. %), E-P-B terpolymer or E-P random copolymer greater than or equal to 70 percent (wt. %), a mixture of glycerol monostearate (GMS) and diethanolamine (DEA) greater than or equal to 0.2 percent, special additive greater than or equal to 0.1 percent, an antiblocking agent greater than or equal to 0.2 percent (wt. %) of synthetic silica or zeolite and an antimicrobial agent greater than or equal to 0.1 percent (wt. %) of Ag + -containing inorganic polymer of linear structure. More preferably, each component of skin layer (A) has a percentage in the following ranges (the total of all components for any specific embodiment would, however, equal 100 percent (wt. %)): polypropylene between 1 and 5 percent (wt. %), E-P-B between 90 and 98 percent (wt. %), a mixture of GMS and DEA between 0.2 and 0.5 percent (wt. %) where the GMS concentration in the mixture may vary from 1% to 99%, and special additive (a mixture of higher fatty acid ester of polyvinyl alcohol or polyether polyol, where respective ratios may vary from 1% to 99%) between 0.1 and 0.5 percent (wt. %) and an antiblocking agent between 0.1 and 0.25 percent (wt. %) of synthetic silica, polymethylmetacrylite or zeolite and an antimicrobial agent between 0.2 and 1.0 percent (wt. %) of Ag + -containing inorganic polymer of linear structure. [0029] Preferably, inner layer (B) has the composition totals: polypropylene greater than or equal to 1 percent (wt. %), E-P-B terpolymer. or E-P random copolymer greater than or equal to 70 percent (wt. %), and a mixture of glycerol monostearate(GMS) and diethanolamine(DEA) greater than or equal to 0.2 percent, special additive greater than or equal to 0.1 percent, More preferably, each component of antifogginner layer (B) has a percentage in the following ranges (the total of all components for any specific embodiment would, however, equal 100 percent (%)): polypropylene between 1 and 5 percent (%), E-P-B between 90 and 98 percent (%), a mixture of glycerol monostearate(GMS) and diethanolamine(DEA) greater than or equal to 0.2 percent, special additive greater than or equal to 0.1 percent, an antiblocking agent greater than or equal to 0.2 percent (wt. %) of synthetic silica. More preferably, each component of antifogging inner layer (B) has a percentage in the following ranges (the total of all components for any specific embodiment would, however, equal 100 percent (%)): polypropylene between 1 and 5 percent (%), E-P-B between 90 and 98 percent (%), a mixture of GMS and DEA between 0.2 and 0.5 percent (wt. %) where the GMS concentration in the mixture may vary from 1% to 99%, and special additive (a mixture of higher fatty acid ester of polyvinyl alcohol or polyether polyol, where respective ratios may vary from 1% to 99%) between 0.1 and 0.5 percent (%) and an antiblocking agent between 0.1 and 0.25 percent (wt. %) of synthetic silica or zeolite. This inner layer does not have antimicrobial agent. [0030] In a preferred embodiment, second inner Layer D has a preferred composition of 100 percent (%) E-P-B terpolymer or E-P random copolymer Further, outer layer E has the same preferred and more preferred compositions as either of layers A or D. [0031] Finally, a preferred embodiment of core layer C comprises the following compositions (the total of all components for any specific embodiment would, however, equal 100 percent (%)): polypropylene greater than or equal to 95 percent (wt %), a mixture of GMS and DEA , would be greater than or equal to 0.2 percent (wt %) and special additive equal or greater than 0.1 percent (wt %) More preferably, each component of layer C has a percentage in the following ranges: polypropylene between 97.5 and 99.5 percent (wt. %), a mixture of GMS and DEA between 0.2 and 0.5 percent (wt. %) where the GMS concentration in the mixture may vary from 1% to 99%, and special additive (a mixture of higher fatty acid ester of polyvinyl alcohol or polyether polyol, where respective ratios may vary from 1% to 99%) between 0.1 and 0.5 percent (wt. %) [0032] Antifogging and Ag + -containing antimicrobial biaxially oriented polypropylene (BOPP) films, can be prepared by using the tandem extruder system with two extruders supplied with two, three or four satellite co-extruders, flat die, chill roll, corona discharge (onto the skin layer or, alternatively, both the skin layer and the outer layer) and recycling line as well as the mono- and semi-oriented cast film technology with temperature controlled mold. After mono- and biaxially stretching (4-7 times at 105-140° C. in the machine direction, MD and 7-11 times at 150-190° C. in the transverse direction, TD) and air corona discharged of one outer surface in the given conditions. Preferably, the antifogging and antimicrobial films have the following characteristics: specific density of 0.91 g/cm 3 , low haze around 1.5% (+−0.2), high gloss greater than or equal to 95%, heat sealability around 120° C., excellent dyne level retention (preferably equal or greater than 40 dynes/cm) for good printability and antifogging characteristics, excellent antifogging properties (rated ‘E’ according to ICI's cold fog test method, which means that the antifogging surface of the film is almost free of big water droplets which makes it invisible) and excellent antimicrobial activity (99.9%) toward various microorganisms, especially and preferably against three common bacteria Staphylococcus aureus, Escherichia. coli and Salmonella enteritidis. [0033] One the other hand, this invention also provides longer shelf life for the freshcut and pre-packed vegetables, salads, fruits and like, due to the high biological activity of the antimicrobial agent which prevents the certain bacteria's growth. Another advantage of the present invention is the easy processability of the antimicrobial agent whose processing conditions are within the processing windows of the ingredients put in the conventionalal or antifogging BOPP films. In fact this advantage is provided by the high thermal stability of the antimicrobial agent which is >300° C. and which is well above of the operating temperatures of the raw materials present in BOPP or antifogging BOPP films. In other words, under normal processing conditions of BOPP film manufacturing, the said antimicrobial agent does not show chemical degradation or decomposition. [0034] Another important advantage of the present invention is the high degree of antimicrobial performance against the certain bacteria by using only a very low concentration of antimicrobial agent, due to its usage only in the very thin layer(s). This very low concentration of antimicrobial agent in the polymer matrix is preferably 30 times lower, as compared with conventional biocides used in known polymer compositions, due to its only use of a thin antifogging and antimicrobial skin layer (A), preferably between 0.5-1.5 μm. The usage of so low concentration of antimicrobial agents advantageously reduces the cost of the film. [0035] Another aspect of the present invention is the possibility to production of the films in the form of mono-oriented and biaxialy oriented multilayer thin films with similar component and layer compositions by using cast film technology and tandem extruder system technology, respectively. According to the present invention the technological aspects of manufactured process of said films are (1) multilayered and mono-oriented cast film technology and (2) tandem extruder system technology by the fact that tandem extruder system with two main extruders for better homogenity and dispersion of the raw materials, supplied with three satellite co-extruders, recycling line and corona discharge. The process is carried out by three chill-roll or water bath treatments and two step of longitudinal orientation allowing to prepare good homogenized film with matte appearance having improved surface properties and dimensional stability. The skin layer or, alternatively both the skin layer and the outer layer, of biaxially oriented films prepared may be treated in a known manner by flame or more preferably, by electrical corona discharge. The use of said recycling line for film waste forming in the transverse stretching stage allows to lower film cost For example (as a preferred, but not the only embodiment of the process), after coextrusion, an extruded five-layer film is taken off over the corresponding process steps through a chill roll and cooled, and cast film profile is controlled by β-Gauge equipment. The film is subsequently stretched longitudinally at two steps and stretched transversely. After biaxially orientation, the film is set and electrical corona-treated on one or two sides. The following conditions are preferrable: (1) Extrusion: extrusion temperature 170-260° C., first chill roll temperature 10-45° C.; (2) Machine (longitudinal) stretching: stretching roll temperature of first step 105-120° C. and second step 115-140° C., longitudinal stretching ratio 4:1-6:1 for first step and 1:1-1:2 for second step; Transverse stretching: temperature of heat-up zones 150-185° C., temperature of stretching zones 155-185° C., transverse stretching ratio 7.5:1-11:1; Recycling: edges of the biaxially orientated film is recycled and fed to the line again; Setting: setting temperature 165-185° C.; electrical corona discharge (A side only or alternatively both A and E sides, together): voltage 10-25 kV and frequency 1.5-30 kHz. The following preferable conditions for the multilayered mono-oriented (in MD only) antimicrobial films in accordance with cast film technology in detail are selected: (1) extrusion temperature 250° C. by using temperature cotrolled MITSUBISHI type die, (2) chill roll temperature 10° C. (3)film profile is controlled by β-Gauge equipment, (3) the speed of film production line 100 m/min, and (4) level of air corona discharge on A surface of the film is 11 Kw [0036] It is further object of the present invention to widen the field of application of said films useful for the food, medicine and agriculture applications as well as for other general packaging and non-traditional special applications including bioprotection of food contacting materials and food handling areas, medicine devices, agriculture products as well as applications in potential areas like food-storage containers, in oral hygienic products, hospitals and other health institutions to provide hygienic conditions, for preserving drinking water and as a covering for water tanks, etc. [0037] Another aspect of the present invention is to use new systems of additives, i.e—a mixture of GMS and DEA, special additive (a mixture of higher fatty acid ester of polyvinyl alcohol or polyether polyol) as antifogging agents, in combination with the antimicrobial agent Ag + , to create dual-effect polymeric films having both antifogging and antimicrobial properties. [0038] Those additive systems are used with the following compatible polyolefins selected from polypropylene, a propylene-ethylene random copolymer, propylene-butene-1 random copolymer or an ethylene-propylene-butene-1 terpolymer with various compositions, where the last three are used for the heat sealable skin layers. [0039] Advantages of antifogging and antimicrobial films are: (1) Antimicrobial activity against certain bacteria, (2) excellent antifogging properties (3) high antimicrobial performance in comparative low concentration of antimicrobial agent, (4) preservation of antimicrobial activity during the long time of storage. of the polymeric films, even after corona or UV-treatment, (5) low total migration properties with the diluents distillated water, acetic acid, ethyl alcohol, heptane and olive oil as mentioned in the directives of EEC and FDA allowing to use of these films in food packagings, (6) high optical properties (low haze, high sheen), (7)high physical-mechanical properties (8) possibility of use various thermoplastic film-forming polymers in core layer of films, and (9) wide range of conventional and special application fields of invented films. DESCRIPTION [0040] The present invention is an antifogging and antimicrobial film that is a multi-layered, oriented and made from polyolefins (polypropylene (PP), propylene-ethylene random copolymers, ethylene-butylene random copolymers and/or ethylene-propylene-butylene (E-P-B) terpolymers with various contents of E- and B- units). The present invention is useful for food packaging, food-wrapping, agricultural and horticultural applications, or any application where there is any condensation of water vapor on the various surfaces in the form of droplets and effectiveness of certain bacteria needs to be reduced. [0041] The following Examples of the present invention for preparation of multilayered antifogging and antimicrobial films with different composition, properties are illustrated. EXAMPLE 1 [0042] A first example of a multilayer film (A/C/E) having antifogging and antimicrobial properties comprises: (A) 1.0 μm antifogging and antimicrobial skin layer containing 92.25% by weight of said ethylene-propylene-n-butylene-1 terpolymer with given composition ( ethylene [C 2 ]=1.5-4.5%, n-butylene-1 [C 4 ]=3.0-15.0%), 0.25% by weight of zeolite as an antiblocking agent, 6.0% by weight of polypropylene homopolymer, and 1.0% by weight of Ag + —as an antibacterial and antimicrobial agent (in derived from a masterbatch having 20% active agent-Ag.+, in a polypropylene carrier: the active is a “silver containing glass powder”, this has the CAS No: 65997-17-3, EINECS No: 266-046-0, and EPA, Reg No: 73148 Issue date: 1 Sep. 2000), 0.20% by weight of glycerol monostearate and 0.20% by weight of diethanolamine and 0.10% by weight of special additive (a mixture of higher fatty acid ester of polyvinylalcohol or polyether polyol) as antifogging and antistatic agent, (C) 28.0 μm core layer (C) from 99.5% by weight of virgin or marked (5-cholesten-3β-ol as a marking agent) polypropylene homopolymer, 0.20% by weight of glycerol monostearate and 0.20% by weight of diethanolamine and 0.10% by weight of special additive (a mixture of higher fatty acid ester of polyvinylalcohol or polyether polyol) as antifogging and antistatic agent and (E) 1.0 μm outer layer having 99.75% by weight E-P-B terpolymer and 0.25% by weight zeolite. This (E) layer does not exhibit any antifogging or antibacterial property. After biaxially stretching the film (5.5 times at 120° C. in the (longitudinal direction, MD and 8 times at 170° C. in the transverse direction, TD) and electrical corona discharged skin layer (A). Layer (A) has corona treatment in order to accelerate the migration of the antifogging agents and alternatively, outer layer (E) has also corona treatment for further printing purposes. EXAMPLE 2 [0043] A second example of a multilayer film comprises the same thickness structure and composition as in Example 1 with the following changes: the core layer (C) comprises 100% of polypropylene homopolymer. EXAMPLE 3 [0044] A third example of a multilayer film comprises an A/C/E structure but with the following changes in Example 1: the antifogging and antimicrobial skin layer (A) is 1.5. μm thick, the core layer (C) is 27.0 μm thick and the non-antifogging, non-antimicrobial outer layer (E) is 1.5. μm thick. After biaxially stretching, heat setting and corona discharged in the given conditions, that film has antifogging and antibacterial properties on skin layer (A), whereas the outer layer (E) is useful for printing and heat seal applications. EXAMPLE 4 [0045] A fourth example of a multilayer film comprises A/C/E thickness structure and composition as in Example 3 with following changes: Outer layer (E) has also antifogging properties but no antimicrobial property. Thus, outer layer (E) has the same antifogging agents in layer (A) of Example 1 but not Ag + which provides antimicrobial effect. Film with that structure is produced as described above. EXAMPLE 5 [0046] A fifth example of a multilayer film comprises A/C/E thickness structure and composition as in Example 4 with following changes: Outer layer (E) has the antifogging and also antimicrobial properties where each of layers (A) and (E) are 1.5 μm thick and have the chemical composition of the skin layer (A) in Example 1. The core layer (C) is 27.0 μm thick and has the same chemical composition as given in Example 3. This film shows antifogging and antibacterial properties and corona treatment on both sides. EXAMPLE 6 [0047] A sixth example of a multilayer film comprises A/B/C/D/E structure with the following changes in the Example 5: inner layer (B) and second inner layer (D) has the same chemical compositions as skin layer (A) and outer layer (E) where each of the four skin layers is of 0.75 μm thick. This symmetrical five layered composition provides the same excellent antifogging and antibacterial properties on both sides. Furthermore, the structure of this example also avoids the low output capacity of the single satellite extruders which limits the total output of the manufacturing line by giving high total extrusion output. Corona discharge on each side of this film also gives the flexibility of using either side by converters or packers. EXAMPLE 7 [0048] A seventh example of a multilayered film comprises A/B/C/D structure with following changes in Example 6: the second inner layer (D) becomes the outer layer (E), having the same chemical composition of the layers (A) and (B) but with a thickness of 1.5 μm. This film exhibits antifogging and antibacterial properties on both sides. EXAMPLE 8 [0049] A eighth example of a multilayer film comprises A/B/C/D/E structure and chemical composition given in Example 6, except the thickness of the core layer (C) which is 32.0 μm, giving the whole structure 35.0 μm total thickness. EXAMPLE 9 [0050] A ninth example of a multilayer film comprises A/B/C/D/E structure and chemical composition given in Example 6, except the following changes: inner layer (B) and second inner layer (D) do not have the antimicrobial agent Ag + , and the thickness of the core layer (C) which is 22.0 μm, giving the whole structure 25.0 μm total thickness. [0051] Layer compositions of the above mentioned examples were given in Table: 1 and the physical-mechanical properties of those films were given in Table:2. [0052] Analysis of the initial materials used and films prepared was done according to known standard measurement methods. For example: Specific density was determined according to ISO 1183 and/or ASTM D-1505. Melt Flow Index (MFI) was measured according to an ASTM 1238/L at 230° C. and under the load of 21.6 N. Melting point (m.p.) was measured by DSC method, maximum point of the melting curve, at a heating rate of 10° C./min. Vicat softening point was determined according to ASTM D-1525. Izod impact strength was measured according to ISO 180/1A. Tensile strength and elongation at,break were determined according to ASTM D-882. Haze of the film was measured according with ASTM D-1003. Dynamic friction coefficient of the film was determined according to ASTM D-1984. Sheen of the film was measured according to ASTM D-2103, the angle of incidence was set at 45°. Shrinkage of the film was measured according to ASTM D-2104. The test sample was shrunk at 120° C. for a period of 5 minutes. Water vapor transmission of the film was measured according with ASTM E96. Oxygen permeability of the film was measured according with ASTM D-1434. Surface tension of the film, after surface ionization by electrical corona discharge and after storage for 6 months, was measured according to ASTM D-2578. Antifogging property of the film was evaluated using ICI's the “Cold-Fog” test method (ICI publication 90-6E) for food packaging film. [0054] The “Cold-Fog” test results of the films according to the present invention (E1-E9), and known patented and commercial antifogging films are summarized in Table 3. The test method is as follows: put tap water, 200 ml, in a 250 ml beaker and cover the top of the beaker with a sample of the test film; place the beaker in a temperature controlled refrigerator at 4° C. Observe the appearance of the film for a total period of one week. It was shown that the films of the present invention, as compared with known patents and commercial films have superior antifogging appearance and properties. [0055] The test method used to measure the antibacterial properties of the present invention is a viable count method. An inoculum, which is a nutrient broth containing a known number of bacteria (there should be 10 5 -10 6 bacteria in the initial inoculum), is placed directly onto the BOPP film. A piece of standard (not antimicrobial) film is placed over the inoculum to ensure intimate contact between the inoculum and the test film and to prevent the inoculum drying out. The sample is covered with the lid of a petri dish and incubated at 35 deg C. and 90% Relative Humidity (ideal conditions for bacterial growth). After incubation the inoculum is washed off the samples, serially diluted and plated out onto Agar plates. These plates are incubated and counts of the still viable (i.e. bacteria able to reproduce and form visible colonies) are counted. Antibacterial test results of the films of the present invention were given Graph 1-3. [0056] Food contact approval tests of the present invention also had been done. Accordingly, Global Migration tests of the preferred embodiment film examples described herein have been found in compliance with the following regulations: EEC Regulation 90/128/EEC and amendments (up to and including 99/91/EEC) and FDA Section 21 CFR Ch. 1 175.300 and 176.170. Those results were tabulated in Table: 4 [0057] According to the present invention, the technological aspect of manufactured process of said films is distinguished from known processing by using the tandem extruder system with two main extruders supplied with two or three satellite co-extruders, recycling line and corona discharge. Other processes of manufacturing said films are known to those skilled in the art. The process is carried out by three chill-roll treatments and two steps of longitudinal orientation followed by the orientation in the transverse direction allowing the preparation of good homogenized antifogging films with improved surface properties and dimensional stability. One or both surface of biaxially oriented films prepared are treated in a known manner by corona discharge. After extrusion, the extruded film having at least 3 layers is taken off over the corresponding process steps through a chill roll and cooled, and cast film profile is controlled by B-Gauge equipment. The film is subsequently stretched longitudinally in two steps and stretched transversely. After biaxially orientation, the film is thermally set and air corona treated on one or two sides. The following are typical manufacturing conditions in detail: (1) Extrusion: extrusion temperatures 170-260° C., first chill roll temperature 10-45° C.; (2) machine direction (longitudinal) stretching: stretching roll temperature of first step 105-120° C. and second step 115-140° C., longitudinal stretching ratio 4.5:1-6:1 for the first step and 1:1-1:2 for the second step; Transverse stretching: temperature of heat-up zones 150-185° C., temperature of stretching zones 155-185° C., transverse stretching ratio 7.5:1-11:1; Recycling: edges of the biaxially oriented film is recycled and fed into the line again; Heat setting: setting temperature 165-185° C.; Air corona discharge: 11 Kw. [0058] While these descriptions directly describe the above embodiments, it is understood that those skilled in the art may conceive modifications and/or variations to the specific embodiments shown and described herein. Any such modifications or variations that fall within the purview of this description are intended to be included therein as well. It is understood that the description herein is intended to be illustrative only and is not intended to be limitative. Rather, the scope of the invention described herein is limited only by the claims appended hereto. TABLE 1 Layer compositions for antifogging-antibacterial films of the present invention. Layer Compositions Exp A skin layer B inner layer C core layer D second inner layer E outer layer E1 1.0 μm — 28 μm — 1.0 μm (1) PP-4.02% PP-97.5% E(2.5%)-P-B(4.5%) (2 )E(2.5)-P-B(4.5%) GMS-0.20% Terpolymer 99.75% Terpolymer-92.25% DEA-0.20% Zeolite-0.25% E-P R.Copo-2.0% E-P R.Copo-2.0% (3 GMS-0.20% Special Additive-0.10% (4 DEA-0.20% (5 )Special Add.-0.10% Zeolite-0.23% Ag+-1.0% E2 1.0 μm — 28 μm — 1.0 μm PP-4.02% PP-100.00% E(2.5%)-P-B(4.5%) E(2.5)-P-B(4.5%) Terpolymer 99.75% Terpolymer-92.25% Zeolite-0.25% E-P R.Copo-2.0% GMS-0.20% DEA-0.20% Special Additive-0.10% Zeolite-0.23% Ag+-1.0% E3 1.5 μm 27 μm 1.5 μm PP-4.02% PP-97.5% E(2.5%)-P-B(4.5%) E(2.5)-P-B(4.5%) GMS-0.20% Terpolymer 99.75% Terpolymer-92.25% DEA-0.20% Zeolite-0.25% E-P R.Copo-2.0% E-P R.Copo-2.0% GMS-0.20% Special Additive-0.10% DEA-0.20% Special Additive-0.10% Zeolite-0.23% Ag+-1.0% E4 1.5 μm 27 μm 1.5 μm PP-4.02% PP-97.5% E(2.5%)-P-B(4.5%) E(2.5)-P-B(4.5%) GMS-0.20% Terpolymer 97.25% Terpolymer-92.25% DEA-0.20% E-P R.Copo-2.0% GMS-0.20% E-P R.Copo-2.0% DEA-0.20% DEA-0.20% Special Additive-0.10% GMS-0.20% E-P R.Copo-2.0% Special Additive-0.10% Special Additive-0.10% Zeolite-0.25% Zeolite-0.23% Ag+-1.0% E5 1.5 μm 27 μm 1.5 μm PP-4.02% PP-97.5% PP-4.02% E(2.5)-P-B(4.5%) GMS-0.20% E(2.5)-P-B(4.5%) Terpolymer-92.25% DEA-0.20% Terpolymer-92.25 GMS-0.20% E-P R.Copo-2.0% % GMS-0.20% DEA-0.20% Special Additive-0.10% DEA-0.20% E-P R.Copo-2.0% E-P R.Copo-2.0% Special Additive-0.10% Special Additive-0.10% Zeolite-0.23% Zeolite-0.23% Ag+-1.0% Ag+-1.0% E6 0.75 μm 0.75 μm 27.00 μm 0.75 μm 0.75 μm ) PP-4.02% PP-4.02% PP-97.5% PP-4.02% PP-4.02% E(2.5)-P-B(4.5%) E(2.5)-P-B(4.5%) GMS-0.20% E(2.5)-P-B(4.5%) E(2.5)-P-B(4.5%) Terpolymer-92.25% Terpolymer-92.25% DEA-0.20% Terpolymer-92.25% Terpolymer-92.25% ) GMS-0.20% GMS-0.20% Special Additive-0.10% GMS-0.20% GMS-0.20% ) DEA-0.20% DEA-0.20% E-P R.Copo-2.0% DEA-0.20% DEA-0.20% Special Add.-0.10% Special Additive-0.10% Special Additive- Special Additive-0.10% E-P R.Copo-2.0% E-P R.Copo-2.0% 0.10% E-P R.Copo-2.0% Zeolite-0.23% Zeolite-0.23% E-P R.Copo-2.0% Zeolite-0.23% Ag+-1.0% Ag+-1.0% Zeolite-0.23% Ag+-1.0% Ag+-1.0% E7 0.75 μm 0.75 μm 27.00 μm 1.50 μm PP-4.02% PP-4.02% PP-97.5% PP-4.02% E(2.5)-P-B(4.5%) E(2.5)-P-B(4.5%) GMS-0.20% E(2.5)-P-B(4.5%) Terpolymer-92.25% Terpolymer-92.25% DEA-0.20% Terpolymer-92.25% GMS-0.20% GMS-0.20% Special Additive-0.10% GMS-0.20% DEA-0.20% DEA-0.20% E-P R.Copo-2.0% DEA-0.20% Special Additive-0.10% Special Additive-0.10% Special Additive- E-P R.Copo-2.0% E-P R.Copo-2.0% 0.10% Zeolite-0.23% Zeolite-0.23% E-P R.Copo-2.0% Ag+-1.0% Ag+-1.0% Zeolite-0.23% Ag+-1.0% E8 0.75 μm 0.75 μm 32.00 μm 0.75 μm 0.75 μm PP-4.02% PP-4.02% PP-97.5% PP-4.02% PP-4.02% E(2.5)-P-B(4.5%) E(2.5)-P-B(4.5%) GMS-0.20% E(2.5)-P-B(4.5%) E(2.5)-P-B(4.5%) Terpolymer-92.25% Terpolymer-92.25% DEA-0.20% Terpolymer- Terpolymer-92.25% GMS-0.20% GMS-0.20% Special Additive-0.10% 92.25% GMS-0.20% DEA-0.20% DEA-0.20% E-P R.Copo-2.0% GMS-0.20% DEA-0.20% Special Additive- Special Additive-0.10% DEA-0.20% Special Additive-0.10% 0.10% E-P R.Copo-2.0% Special Additive- E-P R.Copo-2.0% E-P R.Copo-2.0% Zeolite-0.23% 0.10% Zeolite-0.23% Zeolite-0.23% Ag+-1.0% E.P R.Copo-2.0% Ag+-1.0% Ag+-1.0% Zeolite-0.23% Ag+-1.0% E9 0.75 μm 0.75 μm 32.00 μm 0.75 μm 0.75 μm PP-4.02% PP-5.02% PP-97.5% PP-5.02% PP-4.02% E(2.5)-P-B(4.5%) E(2.5)-P-B(4.5%) GMS-0.20% E(2.5)-P-B(4.5%) E(2.5)-P-B(4.5%) Terpolymer-92.25% Terpolymer-92.25% DEA-0.20% Terpolymer- Terpolymer-92.25% GMS-0.20% GMS-0.20% E-P R.Copo-2.0% 92.25% GMS-0.20% DEA-0.20% DEA-0.20% Special Additive-0.10% GMS-0.20% DEA-0.20% Special Additive- Special Additive-0.10% DEA-0.20% Special Additive-0.10% 0.10% Zeolite-0.23% Special Additive- E-P R.Copo-2.0% E-P R.Copo-2.0% E-P R.Copo-2.0% 0.10% Zeolite-0.23% Zeolite-0.23% Zeolite-0.23% Ag+-1.0% Ag+-1.0% E-P R.Copo-2.0% (1) Polypropylene homopolymer having MFI 1.8-3.5 gr/10 min, at 230° C., under 2.16 Kg. Load, mp = 164-166° C. (2) E-P-B Terpolymer having MFI 5.0-8.5 gr/10 min, at 230° C., under 2.16 Kg. Load, mp = 130-145° C. (3) GMS: Glycerolmonostearate (4) DEA: Diethanolamine (5) Special Additive: Mixture of higher fatty acid acid esters of polyvinyl alcohol or polyether polyol. [0059] TABLE 2 Physical-mechanical properties of the present invention (E1-E9), patented (A), and commercial (B) antifogging films. Patented Films Properties E1 E2 E3 E4 E5 E6 E7 E8 E9 A* B* Total thickness (μm) 30 30 30 30 30 30 30 35 25 19 31 Thickness of core layer 28 28 27 27 27 27 27 32 22 — — (μm) Yield (m 2 /kg) 36.6 36.6 36.6 36.6 36.6 36.6 36.6 31.4 43.9 — — Specific density (g/cm 3 ) 0.91 0.91 0.91 0.91 0.91 0.91 0.91 0.91 0.91 0.91 0.92 Haze (%) 1.7 1.6 1.7 1.6 1.8 1.9 1.9 1.7 1.5 3.1 9.9 Sheen (gloss), 95.2 96.6 96.8 95.2 96.2 95.8 95.3 96.3 97.4 86.6 66.4 45° (%) Shrinkage, 120° C./5 min (%) In MD 3.0 3.5 3.5 3.0 3.0 3.5. 3.0 3.0 3.2 3.5 4.5 In TD 1.0 0.5 1.0 1.0 0.5 1.0 0.5 1.0 1.0 0.25 2.0 Tensile strength at break 13.8 12.6 14.2 12.9 12.3 14.1 12.8 12.3 12.1 13.5 13.6 (kg/mm 2 ), in MD In TD 26.4 24.2 27.3 24.1 27.4 26.8 28.7 25.4 25.7 30.7 25.1 Elongation break (%) In MD 195 193 196 195 198 197 195 185 193 218 183 In TD 58 57 58 59 56 58 58 55 50 50 65.9 Water vapor transmission 4.1 4.3 4.1 4.2 4.4 4.2 4.2 4.0 5.7 ≧15 — (g/m 2 24 h atm 20° C.) Oxygen permeability 1570 1595 1540 1615 1565 1605 1570 1450 1950 ≧3000 — (cc/m 2 24 h atm 20° C.) Friction coefficient, 0.23 0.28 0.25 0.23 0.25 0.27 0.25 0.25 0.22 0.23 0.22 Film/Film Film/Metal 0.20 0.22 0.18 0.20 0.22 0.20 0.22 0.20 0.25 0.22 0.25 Heat seal initiation 120 120 120 120 120 120 120 120 120 125 125 temperature at g/10 mm (° C.) Antifogging property** E E E E E E E E E D C Surface tension(after 40/ 40/ 40/ 40/ 40/ 40/ 40/ 40/ 40/ 37/39 37/32 storage for 6 months) 40 40 40 40 40 40 40 40 40 (nM/m) *A - U.S. Pat. No. 4,876,146 *B - Commercial film **E (Excellent), D (Good), and C (Poor) in accordance with ICI “Cold-Fog” test method. [0060] TABLE 3 Antifogging properties of the present invention (E1-E9), patented and commercial antifogging films. Antifogging properties of films obtained by ICI “Cold-Fog” test method* Antifogging side(s) Example of the examples of No. the present invention Description Performance Rating Comments E1 A A transparent Excellent E Completely film displaying Transparent no visible water E2 A As in E1 Excellent E As in E1 E3 A As in E1 Excellent E As in E1 E4 A, E As in E1 Excellent E As in E1 E5 A, E As in E1 Excellent E As in E1 E6 A, E As in E1 Excellent E As in E1 E7 A, E As in E1 Excellent E As in E1 E8 A, E As in E1 Excellent E As in E1 E9 A, E As in E1 Excellent E As in E1 Patented — Randomly Good D Discontinuous scattered film of water transparent drops Commercial — A complete Poor C Poor Visibility, layer of large lens effect, transparent dripping drops As described in ICI publication 90-6E entitled “Antifog Evaluations Tests for Agricultural and Food-Packaging Film”: [0061] Agricultural and Food-Packaging Film. Description Performance Rating Comments An opaque layer of small fog Very poor A Zero visibility droplets An opaque or transparent Poor B Zero visibility layer of small fog droplets A complete layer of Poor C Poor visibility large transparent droplets Randomly scattered Good D Discontinous large droplets film of water A transparent film with Excellent E Completely no visible water transparent [0062] TABLE 4 Global migration test results of antifogging and antimicrobial film in accordance with EEC Regulation 90/128/EEC and amendments (up to and including 99/91/EEC) and FDA Section 21 CFR Ch. 1 175.300 and 176.170 Food Simulant Test Conditions Mean Result Olive oil 10 days @ 40° C.    2.0 mg/dm2 Distilled water 10 days @ 40° C.    0.2 mg/dm2 3% w/w Ace. Acid 10 days @ 40° C.    0.1 mg/dm2 10% v/v EtoH 10 days @ 40° C.    0.2 mg/dm2 n-Heptan 30 mins @ 70° F.    0.9 mg/in2 Distilled water 24 hrs @ 120° F. <0.01 mg/in2 10% v/v EtoH 24 hrs @ 120° F. <0.01 mg/in2 [0063]

Antimicrobial and antifogging polymeric films with preferable A/C/E structure useful for the food, medicine and agriculture applications as well as for other general packaging and non-traditional special applications. More preferably, antimicrobial and antifogging films having a A/B/C/E structure. Most preferably, antimicrobial and antifogging films having a A/B/C/D/E structure. A multilayer structure (having at least an skin layer (A) having antifogging and antimicrobial properties/a core layer (C)/an outer layer (E) structure) for semi and biaxially oriented polyolefin based antifogging films having advantageous properties as compared with known and commercial films such as low values of haze, high values of sheen, lower longitudinal and transverse shrinkage, which provides high dimensional stability, and excellent antifogging and antimicrobial properties. Preferably the skin layer (A), having antifogging and antimicrobial properties, is electrical corona or flame treated. Electrical corona or flame treatment of the outer layer (E) may enhance ink anchorage and increase the printability of this layer. Preferably, the films comprise an inner (B) layer between the skin layer (A) and the core layer (C). Preferably, the films may comprise a second inner (D) layer between the outer (E) layer and the core layer (C).

This application is a continuation of U.S. application Ser. No. 10/339,829 which was filed on Jan. 10, 2003 and is still pending. BACKGROUND OF THE INVENTION The present invention relates to vacuum cleaners. More particularly, the present invention relates to stick type vacuum cleaners that employ a dust cup. Upright vacuum cleaners are very well known in the art. One type of upright vacuum cleaner which has become common in the marketplace is a stick type vacuum cleaner in which a dust cup is employed for holding dirt and dust separated from an airstream. In most stick type vacuum cleaners, a nozzle base travels across a bare floor, carpet or other surface being cleaned. Pivotally mounted to a nozzle base is an upright housing portion. Often this is formed as a rigid plastic housing having a socket for accommodating the dust cup. As is well known, a suction source such as a motor and fan assembly is mounted to either the nozzle base or the upright housing of the vacuum cleaner. It is now also known in the art of vacuum cleaners to use cyclonic action to separate particles from a stream of dirt laden air. To this end, the dirt laden air is directed tangentially into the dust cup and flows in a swirling motion in the dust cup. Dirt particles are flung outwardly toward the side wall of the dust cup while air is withdrawn along a longitudinal axis of the dust cup. One known type of stick type vacuum cleaner employing a dust cup with cyclonic airflow utilizes an inverted truncated cone positioned within the dust cup. A baffle extends outwardly from an outer surface of the cone. The baffle cooperates with the cone for directing a stream of dirt laden air in a cyclonic manner about the outer surface wall of the cone. In order to remove dust from the dust laden airstream, a filter is positioned outside the dust cup and mounted thereto. This design is disadvantageous from the standpoint that two different elements are needed to provide the cyclonic airflow and to filter the dirt from the airstream. It would be beneficial to have a design wherein the filter element can be positioned in the dust cup rather than being forced out of the dust cup due to the presence of a structure for generating a cyclonic airflow within the dust cup. Accordingly, it has been deemed desirable to develop a new and improved stick type vacuum cleaner which would overcome the foregoing difficulties and others while providing better and more advantageous overall results. BRIEF SUMMARY OF THE INVENTION In accordance with one aspect of the present invention, an upright vacuum cleaner is provided. More particularly, in accordance with this aspect of the invention, the vacuum cleaner comprises a nozzle base including a main suction opening formed on an underside thereof. An upright housing is hingedly connected to the nozzle base. The housing includes a dirt separation chamber and a dirt receptacle for receiving dirt and dust separated by the dirt separation chamber. A conduit connects the nozzle base to the housing. A suction source is located in one of the housing and the nozzle base and is in fluid communication with the dirt separation chamber. A generally conically shaped filter extends into the dirt separation chamber along a longitudinal axis of the dirt separation chamber. According to another aspect of the present invention, a vacuum cleaner is provided. In connection with this aspect of the invention, the vacuum cleaner comprises a nozzle base and a housing pivotally mounted on the nozzle base. The housing defines a cyclonic airflow chamber for separating contaminants from a suction airstream. The housing further comprises an inlet for the cyclonic airflow chamber and an outlet for the cyclonic airflow chamber. A dirt container is selectively mounted in the housing and defines at least a portion of the cyclonic airflow chamber for receiving and retaining dirt and dust separated from the suction airstream in the cyclonic airflow chamber. An airstream suction source is mounted to one of the housing and the nozzle base. The suction source is in fluid communication with the cyclonic airflow chamber and has an inlet disposed downstream from the cyclonic airflow chamber outlet. A filter assembly is selectively mounted to the dirt container and extends into the dirt container. The filter assembly includes a longitudinal axis and a support member including a handle. The longitudinal axis passes through the handle. According to still another aspect of the present invention, a vacuum cleaner comprises a first housing member comprising a cyclonic airflow chamber adapted for separating entrained dirt and dust from the circulating airstream. A dust cup is releasably mounted to the first housing member. The dust cup, which includes an open first end and a closed second end, holds dirt and dust separated from the cyclonic airflow chamber. A second housing member defines a main suction opening. A first conduit fluidly connects the main suction opening of the second housing member to an inlet of the cyclonic airflow chamber. A generally conically shaped filter assembly is selectively mounted to the dust cup. It extends along a longitudinal axis of the dust cup. An airstream source is mounted to the first housing member and is positioned above the cyclonic airflow chamber. The airstream source is adapted for generating and maintaining an airstream flowing through the cyclonic airflow chamber. In accordance with a further aspect of the present invention, a vacuum cleaner comprises a nozzle section and a housing section connected to the nozzle section and in fluid communication with the nozzle section. A dust cup is selectively mounted to the housing section. The dust cup holds dirt and dust separated from a suction airstream flowing into the housing section. A suction source is in fluid communication with the dust cup. A cyclonic airflow chamber is defined at least partially in the dust cup for separating particulate material entrained in an airstream flowing from the nozzle section towards the suction source. A tapered filter assembly extends into the dust cup for further separating dirt and dust from the suction airstream. In accordance with yet another aspect of the present invention, a vacuum cleaner comprises a housing in communication with a suction opening and including a socket. A dust cup is removably mounted in the housing socket. The dust cup comprises an open first end, a closed second end and a side wall. A filter is selectively mounted to the dust cup first end and extends into the dust cup. A particle separation chamber is defined in the dust cup between an interior wall of the dust cup and the filter for separating particles from an airstream flowing from the suction opening through an inlet located in the dust cup side wall. A suction source is in fluid communication with the dust cup first end. The suction source is located in the housing for generating and maintaining a suction airstream from the suction opening through the filter. Still further benefits and advantages of the present invention will become apparent to those of average skill in the art from a review of the following detailed description of the present invention. BRIEF DESCRIPTION OF THE DRAWINGS The invention may take form in certain parts and arrangements of parts, preferred embodiments of which will be described in detail in this specification and illustrated in the accompanying drawings which form a part hereof and wherein: FIG. 1 is a front elevational view of a vacuum cleaner according to the present invention; FIG. 2 is a side elevational view thereof; FIG. 3 is an enlarged exploded perspective view of a lower portion of the vacuum cleaner of FIG. 1 ; FIG. 4 is a rear perspective view of a dust cup of the vacuum cleaner of FIG. 3 ; FIG. 5 is an exploded perspective view of the dust cup of FIG. 3 from above; FIG. 6 is a top plan view of the dust cup of FIG. 5 ; FIG. 7 is an exploded perspective view of the dust cup of FIG. 3 from below; FIG. 8 is a cross sectional view of the vacuum cleaner of FIG. 2 with an upright housing thereof tilted back for use; FIG. 9 is a cross sectional view through the vacuum cleaner of FIG. 2 along lines 9 - 9 ; FIG. 10 is a bottom plan view of the vacuum cleaner of FIG. 1 ; FIG. 11 is a schematic view of a filter according to another embodiment of the present invention; and, FIG. 12 is a schematic view of a filter according to a third embodiment of the present invention. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS Referring now to the figures, wherein the showings are for purposes of illustrating several preferred embodiments of the invention only and not for purposes of limiting same, FIG. 1 illustrates a stick vac A including a nozzle base 10 having on an underside thereof a suction opening 12 . As best shown in FIG. 10 , also provided on the nozzle base are rollers 14 , located immediately behind the suction opening, and rear wheels 16 . With reference again to FIG. 1 , supported on the nozzle base 10 is a housing 20 . Defined in the housing is a socket 22 ( FIG. 3 ) for selectively accommodating a dust cup 24 . Extending from an upper end of the housing 20 is a handle 26 . Positioned on a distal end of the handle is a hand grip 30 which supports a ring 32 . The ring can be used to, for example, hang the stick vac from a suitable peg or hook mounted on a wall or in a closet or the like since the stick vac is a relatively light weight appliance. With reference now to FIG. 2 , also provided on the hand grip is a first cord hook 34 . This cooperates with a second cord hook 36 located on the housing 20 in order to allow a conventional electrical cord (not illustrated) to be wrapped around the cord hooks when the appliance is not in use. With reference now to FIG. 8 , the nozzle base 10 includes a hollow interior 40 which allows air to flow from the suction opening 12 towards a first conduit 42 . The first conduit includes a pivot area 44 at a first end 46 thereof. The first conduit extends out of the nozzle base 10 and terminates at a second end 48 . The first conduit second end 48 is received in a second conduit 50 extending from the housing 20 . To this end, the second conduit 50 has a first end 52 which slips over the fist conduit second end 48 . The two conduits are secured together by conventional means. The second conduit extends along a portion of the housing 20 and terminates at a curved second end 54 which, as is best illustrated in FIG. 3 , leads to an opening 56 . The opening 56 is located on an interior wall 60 of the housing 20 and is surrounded by an elastomeric gasket 62 . As best seen in FIG. 9 , a tangential, swirling, flow path 64 is thus provided for air entrained dirt which flows from the suction opening 12 through the first and second conduits 42 and 50 and out the opening 56 . An upper section 70 of the housing accommodates a motor/fan assembly 72 . This includes a fan 74 and a motor 76 positioned above the fan. Exhaust air from the fan flows out through exhaust vents 78 provided in the housing 20 . With reference again to FIG. 2 , an on/off switch 80 is located on the housing upper section 70 . Also defined on the housing upper section is a catch 84 as best illustrated in FIG. 3 . With reference now to FIG. 4 , the dust cup 24 includes an open first end 92 and a closed second end 94 as well as a side wall 96 extending between the two ends and an interior wall 97 . An opening 98 is defined in the side wall 96 . Extending from the second end 94 of the dust cup is a stub 100 . As best illustrated in FIG. 8 , the stub 100 protrudes into a pocket 101 defined in the housing 20 in order to allow the dust cup 24 to be pivotally mounted on the housing. In other words, the stub 100 and pocket 101 define a hinge assembly for the dust cup on the housing. With reference now to FIG. 5 , a handle 102 is defined on the side wall 96 of the dust cup adjacent the first end 92 thereof. The handle accommodates a moveable, resilient, latch 104 having a finger grip 106 . As best illustrated in FIG. 8 , when a downward force is exerted on the latch 104 at the finger grip 106 , the latch releases from the catch section 84 on the housing thereby enabling the dust cup upper end to be pulled away from the housing. During this time, the dust cup pivots on the housing via a cooperation of the stub 100 in pocket 101 to provide a hinge function. Selectively positioned in the dust cup 24 is a filter assembly 110 . With reference now to FIG. 7 , the filter assembly includes a first end 112 which is defined by a frame 114 . The frame has a top wall 116 ( FIG. 5 ) with a rim 117 and an outer skirt 118 depending therefrom. Spaced from the outer skirt is an inner skirt 119 . Defined in the top wall 116 is a handle 120 which is in the form of a bridge extending between a pair of apertures 122 and 124 in the top wall. Reenforcing the top wall are a plurality of spaced ribs 126 which extend from the outer skirt 118 to the rim 117 . A channel 127 is defined between the inner and outer skirts 119 , 118 . The filter assembly 110 also has a second end 128 which is defined by an end cap 130 having a bottom wall 131 . With reference again to FIG. 5 , extending from the bottom wall is an inner rim 132 and a skirt 134 defining an annular channel 136 therebetween. A filter 140 extends between the frame 114 and the end cap 130 . The filter can be made of a planar filter medium which can be pleated as at 142 . The filter has a first end 144 which is secured in the channel 127 of the frame 114 and a second end 146 which is secured in the channel 136 defined in the end cap 130 , as may best be seen in FIG. 8 . The filter assembly 110 is frustoconical or tapered in its construction. Moreover, the filter material or medium 140 is also tapered in its shape as can be best seen in FIGS. 5 and 7 . It is believed that the conical shape of the filter improves filtering efficiency, as compared with a right cylindrical shaped filter. This may be due, at least in part, to the increased efficiency of cyclonic airflow around the filter that is provided in the dust cup by the cooperation of the dust cup inner wall 97 and the filter. As is evident from FIG. 8 , an approximately constant distance is defined between the filter and the inner wall of the dust cup so as to enhance the cyclonic flow of air around the filter and, hence, dirt separation. The conical filter shape also allows for an easier emptying of the dust cup 24 and may reduce the rate at which the filter 140 becomes clogged. As is evident from FIG. 8 , at the second or smaller diameter end 128 of the filter assembly, the end cap 130 is secured to the filter element 140 . Similarly, at the first or larger diameter end 112 of the filter assembly, the frame 114 is secured to the filter element 140 . The frame 114 includes the pair of apertures 122 and 124 which communicate with an interior 148 of the filter. With this construction, air must enter through the wall of the filter material 140 into the interior space 148 . In other words, the frame 114 and the bottom wall 130 prevent airflow from entering the interior 148 of the filter without passing through the filter medium 140 . The generally conically shaped filter assembly 110 is mounted in the dust cup 24 such that the frame 114 selectively engages the interior wall 99 of the dust cup via an interference fit between the rim 117 of the top wall and the dust cup interior wall 97 . In this way, the filter assembly 110 is releasably, yet securely, retained in its operative position, even when the dust cup 24 is removed from the vacuum cleaner A. Once this is accomplished, the filter assembly 110 can be removed from the dust cup 24 simply by grasping the handle 120 and pulling upward. Thereafter, the dust cup can be inverted so as to remove the dirt and dust contained therein. Subsequently, the dust cup can be righted again, the filter assembly can be reinstalled and the dust cup returned to the housing 20 . This is accomplished by placing the stub 100 in the pocket 101 and then pivoting the dust cup back into position until the latch 104 engages the catch 84 . The latch will be depressed until a tip of the latch clears the backside of the catch and then resiliently snaps upward to hold the dust cup in position, as illustrated in FIG. 8 . The filter material or medium can be made from a suitable conventional planar thermoplastic material if so desired, so that the filter can be washed. Alternatively, the filter medium can be made from a suitable paper material. The frame 114 and end cap 130 can be made from a suitable conventional thermoplastic material. The filter 140 can be secured to the frame 114 and end cap 130 by conventional means, such as adhesive, sonic welding or the like. In use, as best shown in FIG. 8 , air entrained dirt and dust enter the vacuum cleaner via suction opening 12 . The air stream flows through the hollow interior 40 of the nozzle base 10 and into the first end 46 of the first conduit 42 . The air then flows through the first conduit and into the second conduit 50 . Air flows out of the second conduit at its curved second end 54 . The air is directed into a cyclonic swirling flow in the dust cup 24 via the curved second end 54 of the second conduit. The air impinges upon the filter assembly 110 and swirls around it. Thus, a cyclonic airflow chamber 150 is defined in the dust cup between the filter assembly 110 and the interior wall 97 of the dust cup. Particles in the air stream, such as dirt, dust and the like are removed or separated from the suction airstream in the cyclonic airflow chamber. More specifically, the location and orientation of the inlet opening 56 and the generally cylindrical configuration of the cyclonic airflow chamber 150 causes the suction airstream to follow a swirling or cyclonic path within the chamber, as best shown in FIG. 9 . Dirt and dust are flung outwardly by centrifugal force toward the interior wall 97 of the dust cup 90 . The removed particulate matter such as dirt, dust and the like then falls, via gravity, toward the bottom of the dust cup 24 . It is retained therein until the dust cup is emptied. Air, however, flows radially inward toward an axis 152 of the dust cup and then upward around the bottom cap 130 and then radially inward through the filter medium 140 into the interior space 148 thereof. Air then flows upward again through the apertures 122 and 124 around the handle 120 and into the fan 74 . The suction airstream then flows into the fan 74 and out of the housing 20 via the exhaust vents 78 . Thus, a clean air-type vacuum cleaner is here disclosed. As previously noted, the conical or tapered shape of the filter assembly 110 enhances the removal effect of the cyclonic airflow path. Residual particulate matter, i.e., that which is not removed from the suction airstream as a result of the cyclonic action, normally lighter, smaller particles, are filtered by the filter element or medium 140 as the airflow path passes therethrough. The filter assembly 110 extends along the axis 152 of the dust cup such that the filter assembly is centrally positioned in the dust cup. Also, the axis 152 passes through the handle 120 , as may be evident from FIG. 8 . Thus, the filter assembly 110 is concentrically disposed in the dust cup 24 . The location and orientation of the opening 56 in the housing and the opening 98 in the dust cup will effect the direction of cyclonic airflow. However, it is contemplated that the openings could be located and arranged differently. For example, the direction of cyclonic airflow could be reversed. Thus, the cyclonic airflow direction could be clockwise or counter clockwise depending upon the location and arrangement of the aligned openings 56 and 98 . Also, the location of the dust cup side wall opening 98 could be changed if desired. All such orientations and arrangements are considered within the scope of the present invention. Moreover, those skilled in the art will recognize that the term cyclonic as used herein is not meant to be limited to a particular direction of airflow rotation. Rather, the cyclonic action discussed in the present invention is merely intended to separate a substantial portion of the entrained dirt and dust from the suction airstream and cause such dirt and dust to be deposited in the dust cup 24 . The suction airstream then passes through the filter element or medium 140 , so that residual contaminants are removed, and exits the cyclonic airflow chamber, as well as the dust cup, through the two openings 122 and 124 in the frame 114 . One potential disadvantage of the design illustrated in FIG. 8 is that the same portion of the filter medium 140 is exposed to the airstream entering the dust cup 24 . Over time, the dust particles in the airstream may wear the filter material due to prolonged use of the vacuum cleaner. With reference now to FIG. 11 , one way of addressing this issue is to lengthen the skirt of the frame so that the airflow hits the skirt and not the filter medium. More particularly, FIG. 11 illustrates a filter assembly 160 having a first end 162 which is provided with a frame 164 . Extending away from the frame is a skirt 168 . The skirt has a lower end 170 . A filter medium 180 includes an upper end 182 which is in contact with and secured to the skirt lower end 170 along a securement line 184 . The filter medium also has a lower end 186 which is covered by a bottom cap 188 . In the design illustrated in FIG. 11 , the airflow, as depicted by arrow 190 , entering the dust cup (not shown) contacts the thermoplastic material of the skirt 168 and swirls around the skirt rather than directly contacting the filter medium 180 . The material of the frame 164 is less prone to wear than is the material of the filter medium 180 . While FIG. 11 illustrates a design in which the filter is protected from the incoming airstream, a disadvantage of the design illustrated in FIG. 11 is that the filter itself is somewhat shorter, hence, affording less filtration area. With reference now to FIG. 12 , another design is there illustrated. In this design, a filter assembly 200 includes a first end 202 having a frame 204 . Extending from the frame is a skirt 208 . The skirt has a lower end 210 . A filter medium 220 extends away from the frame 204 . The filter medium has an upper end 222 which is secured via a securement line 224 to an inside periphery of the skirt. The filter medium also has a lower end 226 which is covered by a bottom cap 228 . With the design illustrated in FIG. 12 , the airflow, depicted by arrow 230 , contacts the skirt 208 , but yet the length of the filter medium 220 is not shortened. This is accomplished by extending the filter medium upwardly into the skirt until the upper end 222 of the filter is fastened to the skirt via the securement line 224 . The invention has been described with reference to several embodiments. Obviously, modifications and alterations will occur to others upon reading and understanding the preceding specification. It is intended that the invention be construed as including all such modifications and alterations insofar as they come within the scope of the appended claims, or the equivalents thereof.

An upright vacuum cleaner includes a nozzle base having a main suction opening formed in an underside thereof. A housing is hingedly connected to the nozzle base. The housing includes a dirt separation chamber and a dirt receptacle for receiving dirt and dust separated by the dirt separation chamber. A conduit connects the nozzle base to the housing. A suction source is located in one of the housing and the nozzle base. The suction source is in fluid communication with the dirt separation chamber. A generally conically shaped filter extends into the dirt separation chamber along a longitudinal axis of the dirt separation chamber.

This application is a continuation-in-part of U.S. patent application Ser. No. 06/739,321, filed May 30th, 1985, now abandoned. BACKGROUND OF THE INVENTION This invention relates to a container for eggs. SUMMARY OF THE INVENTION It is an object of the present invention to provide a container for eggs, by the use of which the labour involved, in catering establishments and the like, in extracting eggs from the container in which they are normally supplied, boiling the eggs, transferring the same to egg cups for serving and subsequent washing up of the egg cups may be avoided. It is another object of the invention to provide a method of cooking and subsequently serving an egg, utilising a container according to the invention. According to one aspect of the invention there is provided a container for eggs, defining a plurality of compartments each of a configuration to hold an egg captive, each portion of the container defining a respective single said compartment being connected with the remainder along lines of weakening, perforation or the like, whereby each said portion, with an egg retained in the respective compartment, can be readily detached from the remainder of the container, to form a discrete sub-container for a single egg, adapted to hold said egg captive, at least a portion of said sub-container being readily removable to allow access to the egg therein while the remainder of the sub-container acts as an egg cup. According to another aspect of the invention there is provided a method of cooking and subsequently serving an egg comprising placing the egg, whilst enclosed in and held captive in, a permeable container, in water, cooking the egg by boiling the water, removing the container with the egg still retained therein after a desired period, and serving the egg in said container, with said container acting as an egg cup. BRIEF DESCRIPTION OF THE DRAWINGS An embodiment of the invention is described below with reference to the accompanying drawings in which: FIG. 1 is a side elevation view of a container embodying the invention; FIG. 2 is a plan view of the container; FIG. 3 illustrates the lower portion of a sub-container, acting as an egg cup, with an egg in position, said sub-container forming part of the container of FIGS. 1 and 2; FIG. 4 is a perspective view of a further form of egg package embodying the invention; FIG. 5 is a perspective view of an individual egg container forming part of the package of FIG. 4; FIG. 6 is a fragmentary view, partly in side elevation and partly in vertical section, of part of the package of FIG. 4; FIG. 7 is a plan view from below of the individual container of FIG. 5; FIG. 8 is a detailed sectional view, to an enlarged scale, of part of the container of FIGS. 5 and 7, in section along the line V--V of FIG. 7; FIG. 9 is perspective view of another form of egg package embodying the invention; and FIG. 10 is a perspective view of a variant individual egg container. DESCRIPTION OF THE PREFERRED EMBODIMENTS Referring to FIGS. 1 and 2, a container for six egg comprises an upper portion 10 and a lower portion 11, of sheet plastics material deformed to provide, in generally the same way as in conventional egg containers, six egg-receiving pockets in the lower portion 10 and six egg-receiving pockets in the upper portion 10, each lower pocket, in use, serving, with a respective upper pocket, to define a compartment for an egg. As shown in FIG. 2, lines of perforation or weakening 12, 14, in the upper and lower portions 10, 11, divide each egg-containing compartment off from the adjoining compartments, and thus divide the container into six sub-containers 16, the arrangement being such that a desired sub-container may be separated from the remainder by tearing the upper and lower portions 10 and 11, as one, along the appropriate lines 12, 14. In order that each sub-container shall maintain its integrity when torn off from the remainder of the container, the upper and lower parts thereof are held together by appropriate retaining means, for example, by inter-engaging press-fastening formations, 18, of a type known per se in egg containers, provided at the four corners of each sub-container as viewed in FIG. 2, or possibly provided only at two diagonally opposite corners. Both the upper and the lower portion of each sub-container is preferably provided with perforations or apertures 20 therein. The supply of eggs in egg containers of the form shown in FIGS. 1 and 2 allows the boiling of individual eggs to be effected very simply by tearing, from the container, the respective sub-container containing an egg and placing the sub-container, with the egg therein, in boiling water for the appropriate period of time, (usually four minutes), the apertures 20, and the gaps normally present between the upper and lower parts of the sub-container, allowing the entry of boiling water to cook the egg. After having been boiled for the requisite time, the sub-container, with its egg, may be lifted from the boiling water and served directly, the person who is to eat the egg simply removing the upper part of the sub-container from the lower by pulling off in order to gain access to the egg, the lower part of the sub-container, as shown in FIG. 3, serving as an egg cup. The two portions of the sub-container are, of course, simply discarded after use, so that the necessity for washing up egg cups is avoided. As shown in the drawings, the lower part of each sub-container, which is to serve as the egg cup, is preferably made of greater depth than the upper part, for effective support of the egg during use as an egg cup, and the lower part of each sub-container preferably has its peripheral wall formed with vertically extending flutes or corrugations in order to provide thermal insulation between the egg supported by the lower part of the sub-container and the hand of a person supporting the lower part of the sub-container from the outside. If desired, instead of the whole of the upper part of the sub-container being removed to allow access to the egg, an upper portion of each upper part of each sub-container may be connected with the remainder via a line of weakening or perforation and may have a portion affording a tag which can be grasped to allow said upper portion to be torn off, thereby allowing access to the egg, which is nevertheless held captive within the sub-container. It will be appreciated that since the opening of the six-egg container in the conventional manner by pivoting the upper part relative to the lower part about a flexible "hinge" strip along one side of the container, such a hinge strip will not normally be provided so that the upper and lower parts of the six-egg container will normally be connected only via the cooperating fastenings 18. There is thus not, in principle, any limitation on the number of egg spaces which may be provided in a container, and for catering establishments and the like the use of the container affording 24 egg spaces, (i.e. 6×4) may be preferred. Referring to FIG. 4, a package for eggs comprises a rectangular tray 110 provided with circular apertures in which are fitted the lower parts of respective egg containers 112. As shown in FIGS. 5 to 8, each egg container 112 comprises an upper part 114 and a lower part 116 each formed by drawing or moulding sheet plastics material into the form of a cup having around the rim thereof a planar flange 118, 119 respectively having, in plan, the form of a square with rounded corners, said flange lying substantially in the plane of the stock sheet material from which the respective part of the container is formed. The lower part 116 of each egg container has a relatively wide base 120 but has, as viewed externally (see FIG. 7) four channels 122 extending obliquely from the side wall of the lower container part to the base 120, these channels affording, on the interior of the container, corresponding ribs affording inclined support surfaces 124 to support an egg. The lower part 116 is thus relatively stable, allowing the container, or the lower part 116 to stand on its own with an egg therein and yet is configured to provide, internally, adequate support for an egg. The upper part 114 of each container may be, as shown, somewhat more rounded as viewed in side elevation as it does not have to be capable of standing in a stable manner on its own, but if preferred it may have substantially the same form as the lower part 116. In the closed condition of the container, the rectangular flanges 118, 119 of the container parts 114 and 116 respectively are superimposed in parallel and closely adjoining relationship by means of snap-fit formations provided in the four corners of the rectangular flanges, each said formation comprising, as shown in FIG. 8, a moulded stud or pimple 126 extending upwardly from the lower flange 119 and snapped resiliently into a complementary aperture 128 formed in the bottom of a depression 130 in the upper flange 118. The peripheral wall of the stud or pimple 126 is flared upwardly and outwardly slightly so as to secure the desired snap fastening action by virtue of the resilience of the plastics sheet material. The lower part 116 has a central aperture 130 in the base thereof, the upper part 114 has a central aperture 132 in its top and these apertures, together with the spaces defined between the flanges 118 and 119 (see FIG. 8) ensure that when the individual container, with an egg therein, is dropped into boiling water, the air within the container may escape and boiling water may have access to the egg to cook the same. The tray 110 is also formed from sheet plastics material by moulding or drawing and simply comprises a flat rectangular top wall surrounded by a downwardly depending side wall 140, the top wall being formed with six circular apertures, each surrounded by a downwardly depending circumferential wall which forms a collar or sleeve which embraces, as a frictional fit, the side wall of the lower part 116 of a container 112 pushed therein from above. In the embodiment of FIG. 4, the holes in the tray 110 are arranged in a three by two array and the pitch between adjacent holes is such that when containers 112 are inserted in the holes the free edges of the sides of the flanges 118, 119 of adjacent containers lie closely adjacent to one another. In the arrangement of FIG. 9, the tray 110' has six holes therein arranged in a single row to support six egg containers 112, likewise arranged in a single row. The tray 110' is provided at one end with a smaller hole 142 by which, for example, the tray, with the eggs therein, may be suspended from a convenient support. The supply of eggs in egg containers of the form shown in FIGS. 4 to 9 allows the boiling of individual eggs to be effected very simply by removing, from the tray 110, 110', the respective container 112 containing an egg and placing the container 112, with the egg therein, in boiling water for the appropriate period of time, (usually four minutes), the apertures 130, 132, and the gaps present between the upper and lower flanges 118, 119 allowing the entry of boiling water to cook the egg. After having been boiled for the requisite time, the container, with its egg, may be lifted from the boiling water and served directly, the person who is to eat the egg simply removing the upper part 112 of the container from the lower by pulling off in order to gain access to the egg, the lower part 116 of the container serving as an egg cup. The two portions of the container 112 are, of course, simply discarded after use, so that the necessity for washing up egg cups is avoided. FIG. 10 shows an alternative form of the egg container 112, in which the upper container part, indicated at 114a is connected with the lower container part, indicated at 116a by a bayonet-type connection formed by providing the lower container part with a rolled-over peripheral flange 119a having a lower portion extending outwardly from the upper end of the body of the lower container part, and an upper portion extending inwardly again from the outer part of said lower portion to afford a radially inwardly open annular channel 121a defined between said upper and lower portions of said peripheral flange, within which channel 121a are engaged lugs 118a extending outwardly from the lower edge of the upper container part 114a. The upper portion of the rolled-over flange 119a has slots 12a cut in its inner edge portion which correspond in angular spacing and extent to the lugs 118a, so that the upper and lower container parts can be brought together or separated with the lugs 118a aligned with said slots 120a and, once brought together, can be secured by twisting the upper part relative to the lower part to move the lugs out of register with the slots. Preferably the plastics material of the containers of FIGS. 1 to 3 or of the containers 112 of FIGS. 4 to 10 is adapted to undergo a readily visible colour change after it has been subjected to boiling water for the time required to cook an egg, e.g. for four minutes, so that the cook can determine, simply from observing the colour of the plastics container being boiled, with its egg, whether or not the egg has been fully cooked. The desired colour change may be afforded either by incorporating a suitable pigment or chemical in the plastics material itself, or, alternatively, each container may have affixed thereto a piece of an appropriate indicator material. In order to facilitate the cracking of an uncooked egg the inner walls of both the upper and lower portions of each container may be adapted to grip the egg. In this way the two portions of the container provide a reinforcement of the egg shell so that a firm hold can be taken of the egg without risk of crushing the egg during the cracking operation, For example the inner walls could be sculptured to provide a friction grip on the egg that is sufficient to provide a purchase on the egg, but is sufficiently weak to enable the upper portion of the container to be easily removed as described above so that the lower portion can function as an egg cup. When an egg-cracking operation has been completed, the yoke and white emptied from the shell, both the shell segments and the two portions of the container can be discarded. Egg containers as described with reference to the drawings may be utilised in the home, but would be particularly suitable for establishments in which large scale catering is undertaken,for example in restaurants and hotels, as well as in hospitals, for example, old-age pensioners and children's hospitals, in schools, etc. It will be appreciated that, if desired, all of the eggs in a six-egg container or a larger multiple egg-container may be cooked simultaneously in a suitable vessel, with, in the case of FIGS. 1 to 3, the sub-containers being torn off after cooling but just prior to serving, or, in the case of FIGS. 4 to 10, the containers 112 being removed from the trays 110, 110' after cooking but just prior to serving. The packages described may thus, for example, be useful in the preparation of picnics etc. It will be understood that the above description of the present invention is susceptible to various modifications, changes and adaptations, and the same are intended to be comprehended within the meaning and range of equivalents of the appended claims.

A container for eggs defines a plurality of compartments each of a configuration to hold an egg captive, each portion of the container defining a respective single said compartment being connected with the remainder along lines of weakening, perforation or the like, whereby each said portion, with an egg retained in the respective compartment can be readily detached from the remainder of the container, to form a discrete sub-container for a single egg which is adapted to hold said egg captive, at least a portion of said sub-container being readily removable to allow access to the egg therein while the remainder of the sub-container acts as an egg cup. The invention also provides a method of cooling and subsequently serving an egg whilst the egg is enclosed in a permeable container.

BACKGROUND OF THE INVENTION In a number of industries, furniture being one of them, there has been and there is still a need for a time delay switch that is responsive to a physical act, such as pushing a piece of wood stock into a predetermined position to start a desired operation (a sawing or wood carving operation for example) and after a predetermined length of time (time for the operation to be finished) the sawing or wood carving mechanism is automatically de-energized and this de-energized state is maintained while the stock is withdrawn from the working area and the switch repositioned for a subsequent like operation. Only after another piece of stock is again inserted into the same predetermined position is the switch activated to repeat the operation. Accuracy and variability of the predetermined time of operation (or non-operation if desired) is highly desirable in such switch and the means that is used to fix and to change the "on" or "off" time interval should be reliable, rugged, easily adjustable and precise. Repetitive actions, like that described above, require a switch that can be either stationary or mobile and one that can be "re-set" in preparation for a subsequent action after having been rendered to a predetermined state ("on" or "off") for a given time. Resetting must take place in response to the removal of the work stock from the place where it is being subjected to active work and while the switch itself is and is maintained in a predetermined (for example "off") state. It is towards the solution of these problems that the present invention is directed. BRIEF DESCRIPTION OF THE INVENTION The present invention is a time delay switch comprising a housing, an actuating means partially in and partially outside of the housing, a switch capable of being rendered to a first or a second state, biasing means and a compressible/expansional pneumatic means disposed in the housing for rendering the switch from a first to a second state. The actuating means is movable within the housing from a first to a second position from a portion of it outside of the housing. The switch is affixed to the actuating means and the biasing means is connected to the actuating means portion inside of the housing and to the housing itself so that the actuating means is biased towards its first position. The pneumatic means is aligned with the switch and compresses same against the pneumatic means causing the switch to be in its first state when (a) the actuating means is in its first position and when the pneumatic means is in its compressed state; and, (b) when the actuating means is in its second position and the pneumatic means is in its fully expanded state. The penumatic means has a variable rate of expansion from its compressed to its expanded state which is less than the rate of travel of the actuating means can be operated, thereby creating a time interval between the compressed and expanded state whereby the switch is not in contact with said pneumatic means and is thus in its second state. Also included in the pneumatic means is a means for varying the rate of travel from the compressed to the expanded state and this means extends from the pneumatic means inside of the housing to a position on the outside, thereby permitting easy access and permitting an operator a means to adjust the time the switch is to be in its second state. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is a cutaway plan view of one embodiment of the invention. FIG. 2 is a cutaway plan view of another embodiment of the invention. FIGS. 3 and 4 are fragmentary exploded views of the relative positions of pneumatic means 16 and switch 3 of FIGS. 1 and 2 when the pneumatic means 16 is in its fully compressed state. FIGS. 5 and 6 are fragmentary exploded views of the relative positions of pneumatic means 16 and switch 3 of FIGS. 1 and 2 when the pneumatic means 16 is in its fully expanded state. FIG. 7 is a cross sectional front elevation view of an exemplary pneumatic means. DETAILED DESCRIPTION OF THE INVENTION Elements 1 and 35 (FIGS. 1 and 2 respectively) identify first and second embodiments of the present invention. According to FIG. 1 and element 1, housing 2 having sidewalls 41 is shown containing a pneumatic means 16, switch 3 and an "L" shaped actuating means made up of elements 6, 7 and 8, pivotally mounted arm 9, bar 36 and outside member 12. Also included in the housing is biasing means 11, normally a spring. Referring briefly to the pneumatic means 16 as shown in FIG. 7, such device is one that is known in the prior art and can be purchased from the Airpot Corporation of 27 Lois Street, Newark, Conn. 06851 sold underneath the trademark "Airpot." It is composed of a housing 31 which is a precision bore-ultra low friction cylinder in which there is disposed a connecting rod 28 connected to a selectively matched graphite carbon piston 30. The connecting rod is connected to ball joint 29. On one terminal free edge of cylinder 31 is an infinitely adjustable orifice 32, in which there is disposed a mounting stud 33 in which there is disposed screw 18. The space to the left of screw 18 creates vent path 27 through which air may flow into the cylinder 37 to allow the connecting rod to travel upwardly by means of spring 19. Attached to spring 19 is head 17 which is adapted to connect with, depress and release activating switch member 4, as will be hereinafter described. Piston 30 creates two chambers in housing 31, see elements 37 and 40, chamber 37 is in communication with the air outside by means of vent 27. The operation of pneumatic means 16 is believed to be evident from its physical makeup. When element 17 is compressed against spring 19, air is forced out of chamber 37 through vent 27 and spring 19 is compressed. When the compressive forces are released, air from the outside flows through vent path 27 into chamber 37 thereby allowing the force of spring 19 to force head member 17 in the bias direction of the spring. Obviously, the rate of travel of connecting rod 28, spring 19 and head member 17 is a function of the size of vent path 27, which controls the rate of flow of air from the outside into chamber 37. The rate of air transfer is determined by the diametric clearance between piston 30 and cylinder wall 31 and by the vent 27 setting. As piston 30 moves in response to an exerted force, there will be a change in volume and pressure in pneumatic means 16, causing ambient air to enter or leave cylinder 37. By simple adjustment of vent 27, the rate of air flow is controlled to provide an exact degree of head 17 rate of travel. See U.S. Pat. No. 3,175,646. Returning now to element 1 of FIG. 1, pneumatic means 16 is shown disposed in the right-hand side of the housing interior, with adjusting means 18 (screw) disposed on the outside of the housing, so that the rate of flow of air into chamber 37 can be controlled from the outside and thereby control the rate of travel of connecting rod 28, spring 19 and head 17, the importance of which will be hereinafter described. Switch 3 is affixed to arm 7. At the juncture where arms 6 and 7 meet there is affixing means 8 which pivotally attaches arms 6 and 7 to housing 2. Electrically connected to switch 3 are conductor means 5 which are circumscribed by in jacket 14 in a manner well known to the art. Switch 3 can be of any variety, capable of a first and second state, e.g., "on" or "off." Depending on the desired function, depression of activating means 4 can either turn the switch 3 on or it can turn the switch 3 off. An actuating means is shown made up of the "L" shaped member (arms 6 and 7), arm 9, bar 36 and outside means 12. Outside means 12 is pivotally attached to bar 36 and bar 36, at its extreme lefthand terminal portion, is pivotally attached by pin 42 to housing 2. Bar 36 is adapted to abut against arm 9, which is pivotally attached to arm 6 of the "L" shaped means. To a terminal portion of arm 6 is fixed stud 15 onto which is threaded spring 11 and this spring is of such a length that it abuts against housing shoulder 23 at point 24. On the outside of housing 2 is outside means 12, which has a terminal free edge 26 adapted to abut against shoulder 13, thereby limiting its travel and thus the inwardly travel of bar 36, arms 9 and 6 and 7. Effectively, shoulder 13 is a stop means that delimits the inward travel path of outside means 12 and the balance of the activating means attached or abutting thereto or therewith. Switch 3 of FIG. 1 is shown in the second position, e.g., the "on" position. For the sake of description only, the fully extended position of activating means 4 of switch 3 will be hereinafter used and defined as the second state and the fully depressed condition of activating means 4 of switch 3 being the first state, "on" and "off" positions respectively. It is to be realized, however, that the use of the terms "first" and "second" positions is arbitrary and such usage can mean either "on" or "off" relating to the electrical condition of switch 3, depending on the ultimate use of the switch. Referring now to FIG. 2, shown by element 35 is another embodiment of the present invention. This embodiment also has a housing 2 including sidewall 41, in which there is disposed biasing means 11 (usually a spring) terminated at one end to the housing by element 21 (a screw) and on its other end to stop means 20 by bolt 22. Stop means 20 is an integral part of arm 38, which performs a function like arm 6 and 7 of the embodiment 1. Stop means 20 is adapted to abut against a shoulder 25, which delimits its outward travel path from its fully extended to its fully bias state, which will be described hereafter. Switch 3, having activating means 4, is attached to stop means 20 and contains the usual electrical conductors 5 circumscribed by a jacket 14. Switch 3 of FIG. 2 operates in the same way as switch 3 in FIG. 1. Disposed in cavity 34 is a pneumatic means 16 (not fully shown) having a head portion 17. It is constructed in the same manner as element 16 of FIG. 1, except that head 17 is cylinder-like rather than "mushroom" shaped. Disposed on the outside of housing 2 is outside means 12, attached to actuating means 38 by bolt means 39. Outside means 12 contains a terminal free edge 39 which is adapted to abut against shoulder 13 of the housing thereby delimiting the inwardly travel path of outside means 12 and actuating means 38. Turning now to FIGS. 3 and 4, there is shown in these figures switch 3 and activating means 4, in contact with a head 17 of pneumatic means 16, with the spring 19 of the pneumatic means 16 fully compressed. The only difference between elements of FIGS. 3 and 4 is that head 17 of FIG. 3 is "mushroom shaped" (the head 17 of FIG. 4 is cylindrically shaped) and the placement of spring 19. It will be noted that when the pneumatic means 16 is fully compressed, activating means 4 of switch 3 is fully depressed, thereby rendering switch 3 in its first state (e.g., off), which may be either on or off depending on its desired operation or non-operation. A fully compressed state of pneumatic means 16 (shown in FIGS. 3 and 4) is the normal state of the switch; i.e., when there is no pressure or force acting upon outside means 12. When biasing means 11 of element 35 of FIG. 2, is in its fullest expanded state or as with the case of element 1 of FIG. 1 it is in its compressed state, actuating means 38 of FIG. 2 the "L" shaped arm of FIG. 1 (arms 6 and 7) along with arm 9 and bar 36, extend outside means 12 to a position so that it is spaced apart from shoulder 13. Referring to FIGS. 5 and 6, the same elements are shown here as in FIGS. 3 and 4, respectively, the only difference there between being that the pneumatic means 16 in FIGS. 5 and 6 is in its fully extended state, rather than its fully compressed state as was the case in FIGS. 3 and 4. In this state, biasing means 11 of element 1 is compressed, biasing means 11 of element 35 is expanded and outside means 12 in both embodiments abuts against shoulder 13. It will be noted that actuating means 4 of switch 3 in FIGS. 5 and 6 are in the same position (first state) when the pneumatic means is in its fully extended state as was the case when the pneumatic means was in its fully compressed state. Compare FIGS. 3 and 4 with FIGS. 5 and 6. With respect to the aforementioned description of pneumatic means 16, switch 3 and activating means 4 as per FIGS. 3, 4, 5 and 6, a comparison of the state of these elements with like elements of FIGS. 1 and 2 is in order. In FIG. 1, switch 3 and activating means 4 are in the second state and while they are in such second state, it will be noted that the pneumatic means 16 (specifically head 17) is in transition from the fully compressed state as shown in FIGS. 3 and 4 to the ultimately fully extended state (shown in FIGS. 5 and 6). For example, in FIG. 2, activating means or arm 38 can be moved inwardly by outside means 12 from its position where element 12 is spaced apart from shoulder 13 and stop 20 abuts against shoulder 25 to where element 12 abuts against shoulder means 13. In other words, means 38 has moved from its fully biased position to its fully depressed inward position. The rate of travel of acuating means 38 and switch 3 attached thereto, from the fully biased position (stop means 20 abutting against shoulder 25) to the fully inwardly extended position (stop 26 abutting shoulder 13) is greater than the rate of travel of head 17 from its compressed state (FIGS. 3 and 4) to its fully extended state (FIGS. 5 and 6) and thus reaches its fully extended state prior to head 17 coming in contact with activating means 4 of switch 3 and rendering it to the second position or state as shown by FIGS. 5 and 6. The time interval it takes head 17 of pneumatic means 16 to travel from its fully compressed (see FIGS. 4 and 5) to its fully extended (FIGS. 5 and 6) position is that time interval during which switch 3 is in its second state. The rate of travel of head 17, as previously described, is a function of the amount of air going through vent path 27 through orifice 32 into chamber 37 and is variably controlled by screwing screw 18 member inwardly or outwardly so as to increase or decrease the rate of air inflow to lengthen or shorten the time interval it takes head 17 to travel from its fully compressed to its fully extended state, namely, the state as shown by FIG. 4 through the state as shown by FIG. 6. Switch 1 of FIG. 1 works on the same principle as that described for element 35 of FIG. 2, the only difference between embodiments 1 and 35 being the "L" shaped arm (elements 6 and 7) baising means 11, bar 9 pivotally attached to element 6, bar 36 abutting against arm 9, outside means 12 pivotally attached to bar 36 and the position of spring 19 on pneumatic means 16. Elements 12, 36, 9, 6 and 7 of FIG. 1 perform the same function relative to the switch and the pneumatic means as elements 12 and 38 perform for the embodiment 35 of FIG. 2. Inward movement of outside means 12 (its free edge or surface 26 is spaced apart from shoulder 13 in its normal position) pushes against arm 9, which pushes or causes "L" shaped arms 6 and 7 to rotate to the left (counterclockwise) thereby compressing biasing means 11 and removing switch 3 and activating means 4 from contact with head 17 of pneumatic means 16. Upon release of inward pressure of means 12 the reverse takes place. As was described for FIG. 2, the rate of travel of penumatic means 16 from its fully compressed (FIG. 3) to its fully extended (FIG. 5) position determines the time interval that switch 3 in its second state. By rotating screw means 18, this time interval can be increased or decreased as desired by increasing the vent path space through which air can inflow into chamber 37. The dotted lines of FIG. 1, both for the pneumatic means 16 and elements 12, 36, 9, 6 and 7 depict the bias position that these elements would be in in its normal state and the solid lines indicate the position that these elements would be in once inward force has been applied to number 12 sufficient to overcome the biasing means. Operation of the embodiments shown by elements 1 and 35 in FIGS. 1 and 2 are as follows: For the sake of simplicity the structure of element 35 of FIG. 2 will be used; however, it should be kept in mind that elements 12 and 38 of FIG. 2 have the same function relative to the movement of switch 3 and the compression of pneumatic means 16 (head 17) as elements 12, 36, 9, 6 and 7 of element 1 of FIG. 1. The normal state of switch 3, activating means 4 and pneumatic means 16 are as shown by FIG. 4 with elements 38 and 12 are in their fully bias state; i.e., outside means free edge 26 is spaced apart from shoulder 13. Switch 3 in its first state; also, stop means 20 is in abutment with shoulder means 25. On depression of element 12, it moves until its free edge or surface 26 is in abutment with shoulder 13. The relative position of switch means 3 and its activating means 4, visa via the pneumatic means 16 is thus rendered to that instantaneous state as shown in FIG. 2, element 35, this state showing activating means 4 of switch 3 spaced apart from and not in contact with head 17. In such a state, switch 3 is in its second state. Once the relative positions of switch 3 and activating means 4 and head 17 are rendered from that state shown in FIG. 4 to the position shown in FIG. 2, head 17 begins a controlled and much slower rate of travel (slower than the rate of travel) than elements 12 and 38. It travels towards its final position (FIG. 6) and its rate of travel being a function of and dependent upon air from the outside of pneumatic means 16 flowing through vent path 27 into chamber 37, thereby allowing spring means 19 to push head 17 towards activating means 4 of switch 3, while outside air flows into chamber 37. Once head 17 comes in contact with and depresses activating means 4, switch 3 goes from its second to its first state, as shown by FIG. 6. When the inward force (see the inwardly denoted dotted arrow) on element 12 is relieved, biasing means 11 forces arm 36--in the case of FIG. 1 arms 6, 7, 9 and member 36--to return to its normal state. During this return, it will be noted that activating means 4 of switch 3 is still depressed because it is still in contact with element 17. As a result of this contact, switch 3 remains in its first and normal state during the "resetting" or return of switch 3 and pneumatic means 16 to that state shown in FIGS. 3 and 5 respectively for embodiments of FIGS. 1 and 2. In view of the above description, it can be readily seen that a work piece abutted against element 12 and inwardly thrusted causes switch 3 to go from a first to a second state. During this second state, an operation (or non-operation for that matter) can take place and that operation or non-operation is terminated after a predetermined length of time, namely after head 17 depresses activating means 4 of switch 3 thereby rendering switch 3 act to its initial and first state. When the work piece is removed from contact with element 12, the switch is "reset", i.e., the actuating means (elements 12 and 38 in case of FIG. 2) and (elements 12, 36, 9, 6 and 7 in the case of FIG. 1) will return to its normal state; namely, free edge or surface 26 is spaced apart from shoulder 13. During such return, activating means 4 of switch 3 is depressed because it is in contact with head 17. This resetting causes the switch to be in a state to undergo further repetitive operations.

A time delay switch comprising a housing, an actuating device partially in and partially outside of the housing, a switch capable of being rendered to a first and a second state, biasing device and a compressible/expansional pneumatic device disposed in the housing for rendering the switch from a first to a second state; the actuating device is movable within the housing from a first to a second position from a portion of it outside of the housing; the switch is affixed to the actuating device; the biasing device is connected to the actuating device portion inside of the housing and to the housing itself so that the actuating device is biased towards its first position; and, the pneumatic device is aligned with the switch and compresses same causing the switch to be in its first state when the actuating device is in its first position and when the actuating device is in its second position and the pneumatic device is in its fully expanded state, the pneumatic device has a rate of expansion from its compressed to its expanded state less than the actuating device and can be operated, thereby creating a time interval between the compressed and expanded state whereby the switch is not in contact with the pneumatic device and is thus in its second state.

BACKGROUND OF THE INVENTION The present invention relates to a bit counter stage, particularly for memories. More particularly, the invention relates to a bit counter stage with the associated external address loading path for memories and the like. It is known that counters are used in a wide variety of situations, one of the most important being the counting of memory addresses. It is known that a counter is provided by means of a plurality of cascade-connected counter stages, each stage being meant to count one of the bits of, for example, a memory address. The sum of two binary numbers generates a carry value which must propagate along the counter, through the various stages of the counter, in order to obtain a correct sum. The carry calculation time is the fact that limits the operating frequency of a counter. Execution of the carry operation, i.e., its writing time, is the least time-consuming operation; carry generation instead limits and penalizes the operating frequency of said counter. Owing to the need to increase ever more the operating frequency of the counter and therefore to reduce the period of its operation, in conventional counters in which said period is divided evenly between the carry generation step and the carry calculation step the carry generation step, which is the most penalizing one, may not have enough time available for its execution, whereas the carry execution step has an assuredly excessive amount of time available. Moreover, any counter has a loading system which allows to load the initial configuration from outside. It is normally believed that loading management performed by means of an ALE control signal is generally free from reliability problems, but in actual fact there are severe difficulties due to the capability to distribute the current count produced by the counter. In conventional counters, the ALE (address latch enable) address in fact can generate a false count in the counter if the ALE signal is “dirty”, i.e., accidental, or if it is not an actual ALE signal. In this case, the configuration assumed in the counter is destroyed in favor of a configuration which is set externally and the count resumes from the new loading instead of from the current calculation. This should not occur. This situation becomes severe if the counter is directly interfaced with input structures, and in this case the problems involved in preventing possible false updates of the counter become very important. SUMMARY OF THE INVENTION The aim of the present invention is to provide an address counter stage, particularly for memories, in which the cycle of the counter is managed with a time-division mode which minimizes the carry execution interval and maximizes the carry creation and propagation interval. Within the scope of this aim, an object of the present invention is to provide a bit counter stage, particularly for memory addresses, in which the asymmetry of the working cycle of the counter stage allows to raise the operating frequency of said stage. Another object of the present invention is to provide a bit counter stage, particularly for memory addresses, in which the data item in input to the stage is presented to the actual counter stage with a delay technique which is suitable to filter false pulses and is such as to ensure the validity of the operation that must be performed by the counter. Another object of the present invention is to provide a bit counter stage in which the updating steps of each stage of the counter are managed with a technique which guarantees the absence of “chasing” in the various steps and of unwanted acceleration of said steps. Another object of the present invention is to provide a bit counter stage, particularly for memory addresses, in which there is no influence of external stimuli (data) during the counting step of the counter. Another object of the present invention is to provide a bit counter stage, particularly for memory addresses, which can be used with very fast memories and with low supply voltages. Another object of the present invention is to provide a bit counter stage, particularly for memory addresses, which is highly reliable, relatively easy to manufacture and at competitive costs. This aim, these objects and others which will become apparent hereinafter are achieved by a bit counter stage, particularly for memory addresses, comprising: master storage means; slave storage means which are connected to said master storage means; means for enabling the transit of an external address in said master storage means; means for enabling the connection between said slave storage means and said master storage means; means for enabling the connection between said master storage means and said slave storage means; means for calculating the product of said external address and of an input carry signal which arrives from a preceding counter stage; and means for calculating an output carry signal on the basis of said external address and of said input carry signal. BRIEF DESCRIPTION OF THE DRAWINGS Further characteristics and advantages of the present invention will become apparent from the following detailed description of a preferred but not exclusive embodiment of the bit counter stage according to the invention, illustrated only by way of non-limitative example in the accompanying drawings, wherein: FIG. 1 is a block diagram of a counter stage according to the present invention; FIG. 2 is a more detailed block diagram of a counter stage according to the invention; FIG. 3 is a circuit diagram of the first counter stage according to the present invention; FIG. 4 is a circuit diagram of one of the subsequent counter stages, cascade-connected to the first counter stage shown in FIG. 3; FIG. 5 is a circuit diagram of a source of signals for driving the counter stage according to the invention; FIG. 6 plots the timings of the signals of the counter stage according to the present invention; and FIGS. 7 a - 7 c plot the timing of the signals of the counter stage according to the present invention in three different operating modes. DESCRIPTION OF THE PREFERRED EMBODIMENTS With reference to the above figures, and initially to FIG. 1, the counter stage according to the present invention, generally designated by the reference numeral 1 , comprises a bit counter 2 which is suitable to perform the actual counting and receives in input a carry signal CARRY-IN and emits in output a carry signal CARRY-OUT in addition to an address Add. The counter 2 is preceded by storage means 3 which are conveniently constituted for example by a latch which is meant to store an input address 5 which must be transferred to the counter 2 . The address 5 is stored in the latch 3 after passing through a structure of the tristate type 6 , which is driven by a signal ALE-fast (fast address latch enable signal), which enables the loading of the external address 5 into the latch 3 , but the latch 3 does not present the address 5 to the counter 2 itself until there is an actual ALE signal, i.e., an address latch enable signal, which is assuredly a signal of this type. If the ALE signal is not a true ALE signal dedicated to the counter, the external address 5 in fact is not transferred from the latch 3 to the bit counter 2 . Therefore, between the latch 3 and the bit counter 2 there are switching means 7 which are indeed driven by the ALE signal. FIG. 2 is a more detailed block diagram of FIG. 1, in which the latch 3 is shown explicitly and in which the switching means 7 , represented by a pass transistor 8 and by an inverter 9 which receives the ALE signal in input, are likewise shown explicitly. FIG. 3 is a view of a first bit counter stage according to the invention, whereas FIG. 4 is a view of one of the successive stages, cascade-connected to the first stage, which is shown in FIG. 3 . The carry network in the first stage is not present because it coincides with the address itself. FIG. 3 therefore explicitly shows the structure of the counter 2 in addition to the storage means 3 , the switching means 7 and the enable means (tristate structure) 6 . It should be noted that the output carry, designated by CARRY-OUT and CARRY-OUTN, of the first stage and accordingly of all the subsequent stages of the counter constitutes the input carry CARRY-IN of the subsequent stage. The actual counter 2 , shown in detail in FIG. 3 and likewise in FIG. 4, comprises two latch structures, a master latch structure 10 and a slave latch structure 11 , which are mutually connected by first means for enabling the connection between the latch structure 11 and the latch structure 10 , designated by the reference numeral 12 , and of second means for enabling the connection between the latch structure 10 and the latch structure 11 , designated by the reference numeral 13 . The first stage of the counter, shown in FIG. 3, is further provided with an inverter 14 in output to the latch structure 10 . Conveniently, the means for enabling connection between the latch structures 11 and 10 , designated by the reference numerals 12 and 13 , are constituted by pass gates. With reference now to FIG. 4, which illustrates a stage of the counter that follows the first stage, and in which the stage for loading the external address is not shown (even in this stage that follows the first one), clearly there is an external address 5 , enable means 6 , storage means 3 and switching means 7 which are suitable to enable the passage of the address from the storage means 3 to the bit counter 2 . The enable means 12 and 13 are respectively driven by signals S-INC, S-INCn and M-INC, M-INCn, which are generated by a logic network which is shown in FIG. 5 and is constituted by a NOR gate 15 , which receives in input the increment signal INC at one of its inputs, said signal being generated outside the counter stage. A cascade inverter 16 is cascade-connected to the NOR gate 15 and a NAND gate 17 is arranged therein; one input of said NAND gate is constituted by the output of the inverter 16 and the other input of said gate is provided by the signal INC. In output to the NAND gate 17 there is a second inverter 18 , whose output is the signal S-INC, which is fedback to the NOR gate 15 in order to constitute a second input thereof together with the signal INC. The signal M-INC is provided by the output of the NOR gate 15 , the signal M-INCn is provided by the output of the inverter 16 , and the signal S-INCn is provided by the output of the NAND gate 17 . The control signals produced by the network are in pairs and do not overlap, accordingly ensuring the absence of “chasing” effects in the various stages (i.e., “runaway” effects). Going back to FIG. 4, the difference between the inverter stage of FIG. 4 and the inverter stage shown in FIG. 3 is the fact that each counter stage after the first one (such as the one shown in FIG. 4) is provided with a connection between the latch structure 10 and the latch structure 11 which can be provided according to two different paths, one being different from the other due to an inversion produced by the inverter 14 . In practice, each counter stage after the first one has a first path between the master structure 10 , constituted by means for enabling the connection between the master structure 10 and the slave structure 11 which are enabled by the presence or absence of the input carry CARRY-IN and are designated by the reference numeral 20 . Said connection enabling means are conveniently constituted by a pass gate. Likewise, the second path for connection between the master structure 10 and the slave structure 11 , i.e., the path from the output of the inverter 14 to the slave structure 11 , is controlled by second connection enabling means 21 which are also conveniently constituted by a pass gate. The input carry signal CARRY-IN is fed to the N-channel transistor of the connection enabling means 20 , whereas the signal CARRY-INN (i.e., the inverted equivalent of the preceding signal) is fed to the P-channel transistors of the first connection enabling means 20 and to the N-channel transistors of the second connection enabling means 21 . The inverted input carry signal (CARRY-INN) is furthermore fed to carry generating means which are conveniently constituted by a pass gate 22 and by a P-channel transistor 23 , whose gate terminal receives in input the carry signal CARRY-IN. The inverted input carry signal (CARRY-INN) is instead sent to the P-channel transistor of the pass gate 22 , whose N-channel transistor is connected to the gate terminal of the P-channel transistor 23 . In the P-channel transistor 23 , the drain terminal which is connected to means for buffering the output carry (CARRY-OUT), designated by the reference numeral 24 , which emit in output the output carry. The inverted output carry CARRY-OUTN is instead emitted by the pass gate 22 . As shown in FIGS. 3 and 4, the latch structure 11 provides in output the address Add-s, i.e., its content, which must be copied into the master latch 10 after enabling is granted by the means for enabling connection between the latch structure 11 and the latch structure 10 , which are designated by the reference numeral 12 . With reference to the above figures, the operation of the bit counter stage according to the present invention is as follows. In the presence of an address latch enable signal ALE-ext which is presented externally with respect to the counter stage, a signal ALE-fast is generated, likewise externally with respect to the stage, and allows to enable the tristate structure 6 and therefore to make the address ADD-I pass into the storage means 3 , where it is stored temporarily. Then, following the generation of the ALE pulse (which occurs only in case of an authentic request for update from outside), which differs from the address ALE-fast because it is delayed by said address, the switching means 7 are enabled and the address ADD-I previously stored in the storage means 3 is made to pass from said means 3 to the counter 2 , indeed passing through the switching means 7 which are enabled by the signal ALE. This situation is shown in detail in the timing chart of FIG. 6 . With reference now to the diagram of FIG. 3, which illustrates the first stage of a bit counter, it can be seen that once the address ADD-I has passed into the counter 2 , it is in actual fact fed into the master storage means 10 , which are always connected to the slave storage means 11 , by virtue of the means 13 . In this manner, the first stage has a particular configuration, since it has no input carry signal CARRY-IN (and its inverted signal). Then, following the generation of an increment signal INC, which is generated outside the counter stage and is the subject of a patent application in the name of this same Applicant, the means 13 are disabled, setting the signal M-INC low, as shown in FIG. 6, while the signal S-INC is set high, so as to enable the means 12 and therefore allow the connection between the slave storage means 11 and the master storage means 10 . In this manner, the address ADD-I stored in the master storage means 10 no longer passes through the connection enabling means 13 in order to reach the slave storage means 11 . In practice, the connection between the master storage means 10 and the slave storage means 11 occurs in a mutually exclusive manner by way of the connection enabling means 13 and 12 . Therefore, when the connection enabling means 13 are deactivated, i.e., when the path between the master storage means 10 and the slave storage means 11 is interrupted, the path between said slave storage means 11 and the master storage means 10 , across the connection enabling means 12 , is active, allowing to copy from the slave storage means 11 the content that is present therein into the master storage means 10 . In this manner dangerous “runaway” effects of the counter stage are avoided, since otherwise one would have a continuous updating of the address if the paths between the master storage means 10 and the slave storage means 11 were even briefly simultaneously active. The counter stage shown in FIG. 3 (first stage) emits in output both the address ADD and an output carry CARRY-OUT, which constitutes the input carry signal CARRY-IN for a second counter stage which can be as shown in FIG. 4 . All the subsequent counter stages are similar to the one shown in FIG. 4 . With reference therefore to the chart of FIG. 6, it can be seen that the address ADD-I is the current address, whereas the address ADD (output address) is equal to the current address CORR until an increment signal INC occurs and therefore the current address, incremented by means of the carry execution means and stored in the slave storage means 11 , is finally written to the master storage means 10 , when the signal S-INC becomes high. Therefore the output address ADD is updated to the address CORR+ 1 only if the signal S-INC simultaneously becomes high. The updating of the address from CORR to CORR+ 1 is instead performed preemptively and is stored in the slave storage means 11 but is not presented to the master storage means 10 until the signal S-INC occurs. The address that is present in the master storage means 10 therefore has two path options for reaching the slave storage means 11 , depending on whether the input carry signal CARRY-IN is present or not. If the input carry signal CARRY-IN is high, then the connecting path through the pass gate 20 is enabled and therefore the contents of the master storage means 10 pass without inversion into the slave storage means 11 and the output carry is calculated by the transistor 23 by means of the pass gate 22 . Otherwise, the connecting path that occurs after the inverter 14 , through the pass gate 21 , is enabled and therefore the contents of the master storage means 10 undergo inversion through the inverter 14 and then pass through the pass gate 13 and reach the slave storage means 11 . The carry execution means actually provide a multiplication for the input carry CARRY-IN and the data item (address) contained in the master storage means 10 and copy said product to the slave storage means 11 ; both the data item (address) and the input carry signal CARRY-IN contribute to the calculation of the output carry CARRY-OUT, which is then sent to a subsequent bit counter stage, similar to the stage being considered, in order to constitute its input carry signal CARRY-IN. In practice, it has been observed that the external address ADD-IN is fed into the counter exclusively due to the presence of an assured address latch signal ALE; this therefore allows to contain the consumption of the counter stage with respect to conventional solutions in which the counter stage performs the calculation anyway, without realizing whether the address presented externally thereto really is an address that must be updated or is an address which generates a false loading. Therefore, each counter stage is isolated from the outside world by means of the tristate structure 6 and then by means of the switching means 7 , which are driven by an authentic address latch signal ALE. Moreover, the updating of each counter stage is performed at a different time with respect to the copying of the updated address in the master storage means. Moreover, the time dedicated to the calculation of the carry is much longer than the time dedicated to the writing of said carry, and this allows to increase the operating frequency of the counter. During counting, stimuli that are external to the counter are disabled and protection is provided with respect to external loadings, allowing to recognize noise-related or non-authentic signals, thus allowing overwritings of the counter stage only if this is intended. The charts of FIGS. 7 a - 7 c plot different situations in which FIG. 7 a shows the transition from a standby state to an active mode in which the signal ALE is generated with a delay with respect to the signal ALE-ext, and following an inverted read signal RDen the counter (designated in this case by a signal COUNT) updates its own address. The counter stage thus conceived is susceptible of numerous modifications and variations, all of which are within the scope of the inventive concept; all the details may furthermore be replaced with other technically equivalent elements.

A bit counter stage, particularly for memory addresses, including: a master storage circuit; a slave storage circuit which is connected to the master storage circuit; a circuit for enabling the transit of an external address in the master storage circuit; a circuit for enabling the connection between the slave storage circuit and the master storage circuit; a circuit for enabling the connection between the master storage circuit and the slave storage circuit; a circuit for calculating the product of the external address and of an input carry signal which arrives from a preceding counter stage; and a circuit for calculating an output carry signal on the basis of the external address and of the input carry signal.

This is a continuation of application Ser. No. 07/750,868, filed Aug. 28, 1991, now abandoned. BACKGROUND OF THE INVENTION The present invention relates generally to security sensors and to energy-conservative sensors for sensing entry into a monitored area or room. More specifically, the present invention relates to passive sensors for automatically lighting and extinguishing lights when a person enters and leaves a room. Passive infared (PIR) motion sensing is an expanding technology driven by security and energy conservation demands. One typical use of conventional PIR technology is automatic illumination of room lighting when a person enters a room. A timer will automatically extinguish the lights after a predetermined interval unless the PIR detects the person moving in the room. A disadvantage of these motion-only sensors is that the sensor may not detect a presence of person out of a line-of-sight of the PIR but still within the room. In this case, the sensor extinguishes the lights while the person remains in the room. The consequences of this premature loss of light range from inconvenience and annoyance to potential hazard and bodily injury, depending upon the particular room or area monitored. In some instances, a person may merely have to wave or stand to trigger the sensor while in other instances the person way have to move to a monitored part of the room in darkness. These prior art room sensors typically employ the same mechanism for triggering and retriggering. That is, the sensor will illuminate room lights (trigger them) when it detects motion and will reset a timer (retrigger the lights) when it detects motion. SUMMARY OF THE INVENTION The present invention provides apparatus and method for sensing entry into a room or other monitored area. The present invention provides user-determinable preconditions of selected environmental conditions in the monitored area for triggering and retriggering. The triggering and the retriggering are independent from each other and are able to employ different sensors and monitor different parameters. In one preferred embodiment, the sensing apparatus includes a motion sensor, an ambient light sensor, and a sound sensor. The preferred embodiment also includes an actuator and a timer, as well as logic circuitry to test for the desired preconditions for triggering and retriggering. In operation, the preferred embodiment for controlling room illumination monitors for motion in low ambient room light. Upon detecting motion with low light, the logic circuitry triggers the actuator and initiates the timer. To determine when to retrigger, the sensor detects for sound or motion within the room. Without either sound or motion in the room, the timer will expire, extinguishing the lights. After extinguishing the lights, the sensor will wait for its predetermined triggering configuration of the environmental conditions monitored by its sensors. In another aspect of the invention, the sensitivities of the various sensors are adjustable, providing a large range of applications for the present invention. For instance, adjusting a sound sensitivity for a room permits retriggering simply by conversing with another person, or by turning pages of a book or newspaper. It is possible to adjust either mode of operation, triggering or retriggering, so that only a single sensor will monitor the desired environmental condition. Additionally, the present invention permits priorities or particular orderings of selected environmental conditions to trigger or retrigger the actuator in response to the sensors, where the priorities or particular orderings of environmental conditions may be selected, independent of each other. One example sets the sensors so that retriggering results from detecting sound only after first detecting motion. For security areas, in some instances it is desirable to trigger an actuator controlling an alarm or light after detecting motion, a flashlight beam and a sound of forced entry, for instance. In some instances, triggering results from first detecting motion, then forced entry, or vice versa. Proper order of the selected signals results in triggering and retriggering, with each independently selectable. Another embodiment of the present invention employs radio frequency transmissions between the sensors and the actuator, allowing remote switching of desired load. The actuators may operate from power supplies independent from those of the sensors. The present invention provides users with an ability to tailor entry sensors for particular applications. The improved entry sensor enhances convenience and safety of the user, permitting widespread acceptance of illumination and security controls using the present invention. Reference to the remaining portions of the specification and drawings may realize a further understanding of the nature and advantages of the present invention. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is a view of an overall perspective preferred embodiment of the present invention illustrating a configuration for room illumination control in a switch model mounted within a standard wall box; FIG. 2 is a block diagram of a circuit employing the present invention; FIG. 3 is a flow chart illustrating operation of a preferred embodiment of the present invention; and FIGS. 4A-4D show alternate embodiments of the invention illustrating separate use of light, PIR and sound sensors interconnected by a radio frequency (RF) link. FIG. 4A shows an RF-transmitting motion sensor unit including a PIR sensor and light sensor formed to fit in a socket for a spotlight or track lamp. FIG. 4B shows a battery-operated RF-transmitting motion sensor unit including a sound sensor. FIG. 4C shows an RF-receiving motion sensor unit including a sound sensor mounted together with a conventional electrical wall switch. FIG. 4D shows an RF-receiving motion sensor including sound and light sensors mounted together with a conventional electrical wall outlet. FIG. 5 shows a hand-held remote control unit for use with the RF units of FIGS. 4A-4D. DESCRIPTION OF THE PREFERRED EMBODIMENT FIG. 1 is a view of a preferred embodiment of the present invention illustrating a configuration for room illumination control 10 in a switch model mounted within a standard wall box. The illumination control 10 includes a plastic body 20 and a metal mounting plate 31. Conventional mounting of the mounting plate 31 to a wall box 24 with mounting screws 25 through a cover plate 27. Power from a power source, such as household alternating current connects to one line of the illumination control 10 and a second line connects to another line of the illumination control 10. A manual slide switch 37 has different positions (OFF, ON, and AUTO) for the different functions of the illumination control 10. In the OFF position, the illumination control 10 is incapable of activating a load regardless of particular environmental conditions in the room. In the ON position, the illumination control 10 activates the load, again without regard for particular environmental conditions. In the AUTO position, the illumination control 10 begins a passive infared (PIR), light and sound detecting process further explained below with reference to FIG. 3. The illumination control 39 includes a fresnel lens 38 focusing infrared radiation from a monitored area onto a pyroelectric infrared sensor, not shown. The illumination control includes a photoelectric sensor mounted behind transparent cover 32 and a microphone mounted behind holes 33. Sensitivity switches 36, 34 and 35 for each of the three sensors, PIR, photoelectric and microphone respectively adjust a sensitivity of their associated switches. For example, in the preferred embodiment, adjustment sensitivity for switch 35 (audio) permits triggering or retriggering from a range of 10 dB to above 110 dB. FIG. 2 is a block diagram of a sensing circuit 50 employing a preferred embodiment of the present invention. The sensing circuit 50 includes three environmental sensors 52, 54 and 56 for monitoring various environmental conditions, such as motion, light and sound, for example. A logic circuit 60 monitors output signals from each of the sensors. The logic circuit 60 may be implemented in any number of well-known implementations, including microcontroller circuitry or hardwiring. Each sensor 52, 54, 56 has an associated sensitivity adjustment switch 61, 62, 63 used to set threshold levels. The logic circuit 60, responsive to particular configurations of environmental conditions measured by the sensors and a mode of operation, controls an actuator 64. The actuator 64 controls a load 66. The load 66, for example, may be a light, or an alarm. The logic circuit 60 deactivates the actuator 64 responsive to an assertion of a timing signal from a timer 68. A control switch 70 selects whether the sensing circuit is ON, OFF or in AUTOMATIC detection. When ON, the logic circuit 60 causes the actuator 64 to activate the load 66 irrespective of a condition of the output signals from the sensors 52, 54 and 56. When OFF, the actuator 64 deactivates the load 66. In AUTOMATIC, the logic circuit 60 operates as identified in the description relative to FIG. 3. FIG. 3 is a flow chart of the operation of the sensing circuit 50 for implementation of an illumination control. Steps 40 through 48 are process steps implemented by the logic circuit 60 of FIG. 2. Sensor 52 is a PIR sensor, sensor 54 is a photoelectric sensor, and sensor 56 is an audio sensor. The load 66 is a lamp. At step 40, the logic circuit 60 determines whether the control switch 70 is in AUTOMATIC or not. If in AUTOMATIC, the logic circuit 60 determines whether the output signal from the PIR sensor falls within a prespecified range indicating motion within a monitored area, at step 41. The range may optionally include upper and lower bounds. If the PIR sensor does not detect motion, the logic circuit 60 cycles back to the step 40, continually testing for motion within the monitored area. If at step 41, the PIR sensor indicates motion, the logic circuit 60 advances to step 42 to test an ambient light level with the photoelectric sensor. For the preferred embodiment, if the ambient light exceeds a predetermined threshold, the monitored area is sufficiently illuminated so the logic circuit 60 will not trigger the lamp, but return to step 41. However, if at step 42 the ambient light is below the threshold, the logic circuit will proceed to step 43. At step 43, the PIR sensor detected motion, followed by the photoelectric sensor determining that the monitored area was sufficiently dark to warrant further illumination. Thus, at step 43, the logic circuit 43 causes the actuator 64 to activate the lamp. Additionally, the logic circuit 60 resets and starts the timer 68. Passing the logic tests at step 41 and step 42 triggers the illumination control. The timer 68 measures a lapse of a predetermined interval. If the illumination control has not been retriggered before expiry of the timer 68, as indicated by assertion of a timer signal to the logic circuit 60, the logic circuit 60 will cause the actuator 64 to turn the lamp off. The retrigger mode cycles through the steps 44 through 48. In the preferred embodiment, the retriggering begins at step 44 with the audio sensor checking for any sound in the monitored area. Sound falling within an identified range results in the logic circuit 60 resetting the timer 68 at step 48. After resetting, the logic circuit returns to step 44. If the sound level falls outside the identified range, the logic circuit 60 advances to step 45 to test for motion. An output signal from the PIR sensor within a second (retriggering) range will cause the logic circuit 60 to reset the timer 68 at step 48 and return to step 44. Absent sound or motion within the proper ranges, the logic circuit 60 advances to step 46 to test a status of the timer 68. If the timer 68 has not expired, the logic circuit 60 returns to the step 44. However, finding the timer 68 expired, the logic circuit 60 causes the actuator 64 to deactivate the lamp. Thereafter, the logic circuit 60 returns to step 41, waiting for the particular combination of environmental conditions as measured by the particular combination of sensors. As the preferred embodiment of the present invention includes multiple sensors monitoring different environmental conditions, it is contemplated that different environmental conditions include two sensors which both detect motion, but in different areas of a monitored area. The term "configuration of environmental conditions" refers to any permutation or combination of the various parameters measured by the particular sensors employed. Ordering is an important part of a preferred embodiment of the present invention. FIGS. 4A-4D and 5 show alternate preferred embodiments of the present invention illustrating separate use of light, PIR and sound sensors interconnected by a radio frequency (RF) link. The alternate embodiment include a transmitting sensor 102, a receiving sensor 104 and a handheld remote control 106. The transmitting sensor 102 of this preferred invention includes either a combination motion/light sensor 102a (FIG. 4A) or a combination motion/sound sensor 102b (FIG. 4C). The transmitting sensors 102 operate from either conventional a.c. power (such as by a bulb socket) or battery operated. These transmitting sensors 102 are positioned separate from a cooperating receiving sensor 104. The transmitting sensor 102 provides the sensor signals from its sensors to the receiving sensor 104 via radiofrequency, or equivalent such as infrared signalling. The cooperating receiving sensor 104, which is a receiving sound sensor 104a (FIG. 4B) for transmitting sensor 102a or a receiving sound/light sensor 104b (FIG. 4D) for transmitting sensor 102b. In this preferred embodiment, the logic circuitry 60 is included within the receiving sensor 104. The system operates similarly to the sensing system described above, except that the radiofrequency intercommunication permits a wider range of environmental conditions for triggering or retriggering events as the sensors are able to be physically separated. To control a system according to the alternate preferred embodiment, the remote control 106 (FIG. 5) can place the system in any of the ON, OFF or AUTOMATIC states. In conclusion, the present invention provides many advantages over existing devices, including more versatile, convenient and customizable operation. The present invention offers differing triggering and retriggering configurations for various environmental conditions. While the above is a complete description of the preferred embodiments of the invention, various alternatives, modifications, and equivalents may be used. For example, other remote interconnection systems other than radiofrequency can allow communication of the sensor signal. Therefore, the above description should not be taken as limiting the scope of the invention which is defined by the appended claims.

A sensing apparatus and method for use in an illumination control monitoring of a particular area. The sensing apparatus includes selectable and independent triggering and retriggering modes for activating and deactivating a lamp. Three sensors, a PIR, a sound and a light sensor cooperatively interact to cause triggering upon detecting motion in a low ambient light room. Thereafter, retriggering results upon either motion or sound being periodically detected in the room. An absence of motion or sound within a predetermined duration results in deactivation of the light and a return to a pre-triggering configuration.

This Application claims the benefit of U.S. Provisional Application No. 60/565,870, filed Apr. 28, 2004. FIELD OF THE INVENTION The present invention relates to postoperative treatment of cancer using UFT. In particular, the invention relates to improved methods for treating lung cancer by postoperative adjuvant chemotherapy with UFT. BACKGROUND OF THE INVENTION UFT is an oral anticancer agent comprised of tegafur and uracil at a molar ratio of 1 to 4 which has good absorption in the small intestine. (Fujii, S., et al., 1979; U.S. Pat. No. 4,328,229). Tegafur is gradually converted to 5-fluorouracil via the metabolism of liver enzyme P450. Uracil enhances the serum 5-fluorouracil concentration by the competitive inhibition of dihydropyrimidine dehydrogenase, the enzyme responsible for 5-fluorouracil catabolism. (Ikenaka, K., et al., 1979). Oral UFT administration reportedly generates a higher maximum plasma level of 5-fluorouracil than the protracted intravenous injection of 5-flourouracil given in a dose equimolar to the tegafur in UFT. (Ho, D. H., et al., 1998). The response rate of single UFT treatments in patients with advanced stage lung cancer is reported to be 6 to 8 percent. (Shimizu, E., et al., 1986; Keicho, N., et al., 1986). Combination chemotherapy consisting of a daily administration of UFT for two or three weeks, and a bolus injection of cisplatin in advanced non-small cell lung cancer patients yields a response rate of 29 percent to 38 percent and a median survival time of eight to thirteen months. (Ichonose, Y., et al., 1995; Ichinose, Y., et al., 2000; Saito, J., et al., 2001). In two trials for locally advanced non-small cell lung cancer patients, the combination chemotherapy of UFT plus cisplatin with concurrent radiotherapy shows a response rate of 80 percent (Ichinose, Y., et al., 2002) and 94 percent (Ichinose, Y., et al., 1999) and a median survival rate of 16.5 months. (Ichinose, Y., et al., 2002). These results of UFT plus cisplatin chemotherapy regimens are comparable to those of other recently published cisplatin based doublet chemotherapy regimens. (Schiller, J. H., et al., 2002; Vokes, E. E., et al., 2002). Adenocarcinoma, a form of non-small cell lung cancer (NSCLC), accounts for approximately 40% of all cases of lung cancer. It is the most common form of NSCLC and the most common type of lung cancer overall. Concerning adjuvant treatment using UFT, the West Japan Study Group for Lung Cancer Surgery reported that postoperative adjuvant treatment with UFT (400 mg/day for 1 year) in patients with completely resected stage I-III disease prolonged survival significantly longer than observation alone. (Wada, H., et al., 1996). The 5-year survival rate was 64 percent in the UFT group and 49 percent in the control group (P=0.02). In a subgroup analysis, no statistically significant difference in the overall survival of patients with squamous cell carcinoma between the two groups was observed (P=0.24). In contrast, for the patients with adenocarcinoma in the UFT group, most of whom had stage I disease, survival was significantly better than for those in the control group (P=0.009). (Okimoto, N., et al., 1996). Improved methods for extending survival while causing minimal side effects in postoperative lung cancer patients, particularly in patients with adenocarcinoma, are needed. SUMMARY OF THE INVENTION In one aspect, the present invention relates to a significantly improved method for treating lung cancer, preferably non-small cell lung cancer, and more preferably adenocarcinoma, by postoperative adjuvant chemotherapy with UFT orally administered substantially daily to a postoperative lung cancer patient. In a preferred method, about 100–500 mg/m 2 /day of UFT, preferably about 200–300 mg/m 2 /day of UFT, and most preferably an average of about 250 mg/m 2 /day UFT is orally administered to a postoperative lung cancer patient in need thereof for a period of at least about two years. In a preferred embodiment, the present invention relates to treating a patient having pathological stage I, and more preferably pathological stage IB, adenocarcinoma of the lung. In a more preferred embodiment, the present invention relates to treating adenocarcinoma which has been completely resected prior to treatment. In another preferred embodiment, the present invention relates to treating a patient having a primary tumor of a size more than 2 cm, and preferably about 2 to 3 cm, or more, and more preferably wherein the primary tumor is classified as T2 according to the TNM classification (See Mountain, C. F., 1997; Mountain, C. F., et al., 1997). BRIEF DESCRIPTION OF THE FIGURES FIG. 1 shows the survival of all eligible 979 patients (Panel A), 263 patients with T2 Disease (Panel B) and 716 patients with T1 disease (Panel C) assigned to the UFT group and the control group. Error bars represent the 95 percent confidence intervals. The P value was calculated using the stratified logrank test. FIG. 2 shows the interaction of prognostic factors with treatment in survival. Each square estimates the treatment effect, and horizontal lines represent the 95 percent confidence intervals. The diamond corresponds to the 95 percent confidence intervals for whole subjects. P-value of the tumor size is for the interaction with two groups of <=2 and >3 cm. DETAILED DESCRIPTION OF THE INVENTION In accordance with the invention, the oral administration of UFT (a combination of tegafur and uracil at a ratio of 1:4) as adjuvant chemotherapy was shown to prolong the survival of patients with resected adenocarcinoma, among whom most patients had pathological stage I disease. In particular, postoperative adjuvant chemotherapy using UFT (250 mg/m 2 /day) for a period of 2 years was found to yield a significant improvement in the survival of patients with pathological stage I adenocarcinoma of the lung, especially in stage 1B (T2NOMO) (Kato, H., et al., 2003). Patients with completely resected pathological stage I adenocarcinoma of the lung were randomized with stratification according to their pathological T status (T1 versus T2), gender and age to either receive the oral administration of UFT (tegafur 250 mg/m 2 /day) for two years or no treatment. The primary endpoint was overall survival. Patient Characteristics Between January 1994 and March 1997, 999 patients were enrolled in the trial and 501 patients and 498 patients were randomly assigned to receive either no treatment or UFT, respectively. However, seven patients in the UFT group and 13 patients in the control group were found to be ineligible for the following reasons: pathological N1 or M1 disease in seven patients, a histology other than adenocarcinoma in six, no laboratory data at registration in two and other miscellaneous reasons in five. Therefore, the number of all eligible patients was 488 in the control group and 491 in the UFT group. The clinical characteristics of those eligible patients are listed in Table 1. There were no statistically significant differences in the base line characteristics of the patients. All but one patient in each group underwent lobectomy. Adverse Reactions and Compliance Of the 498 patients randomized to the UFT group, 482 patients received an oral administration of UFT. Table 2 lists the incidence of UFT-related adverse reactions. Few severe adverse reactions were associated with UFT administration. There was no grade 4 adverse reaction. In total, 10 (2 percent) of 482 patients developed a grade 3 adverse reaction. The percentage of compliance for UFT administration was calculated based on the number of patients who actually took UFT and the number of patients without recurrence, second cancer or death who were expected to take UFT. The percentage of compliance was 80 percent (95 percent confidence interval: 77 to 84 percent) at 6 months, 74 percent (95 percent confidence interval: 70 to 78 percent) at 12 months, 69 percent (95 percent confidence interval: 65 to 73 percent) at 18 months and 61 percent (95 percent confidence interval: 77 to 84 percent) at 24 months. The main reasons for a discontinuation of UFT administration were as follows: an adverse reaction in 123 patients, patient refusal in 52 and the doctor's judgment in 34. Overall Survival The median follow-up for the surviving patients was 72 months in the UFT group and 73 months in the control group. The numbers of censored patients were 361 in the UFT group and 359 in the control group. At the last follow-up, 65 patients in the UFT group and 89 in the control group had died and the overall survival showed a statistically significant difference based on the stratified logrank test as shown in FIG. 1A . The 5-year survival rate was 88 percent (95 percent confidence interval: 85 to 91 percent) in the UFT group and 85 percent (95 percent confidence interval: 82 to 89 percent) in the control group. When the survival analysis was performed in all 999 randomized patients, the result did not change and the p-value of the difference between the two groups was 0.047. The predetermined covariates were age (<65 years versus >65 years), sex (male versus female), performances status (0 versus 1 plus 2), T status (T1 versus T2) and treatment groups. The covariates were selected according to multivariate analysis using a stepwise procedure under the condition that the p value was less than 0.05. The selected covariates were as follows: age (hazard ratio=2.02, 95 percent confidence interval=1.46 to 2.80 percent; P<0.001), T status (hazard ratio=1.95, 95 percent confidence interval=1.41 to 2.69 percent; P<0.001) and sex (hazard ratio=0.66, 95 percent confidence interval=0.48 to 0.91 percent; P=0.01) and treatment groups (hazard ratio=0.72, 95 percent confidence interval=0.53 to 1.00 percent; P=0.05). An interaction between four prognostic factors listed in FIG. 2 and the treatment was then evaluated. Since the T status is mainly classified by the minimum diameter of the primary tumor, the tumor size was added to the analysis. As shown in FIG. 2 , significant interaction of either T status or the tumor size with the treatment was observed. The patients with T2 disease in the UFT group had a significantly better survival than those in the control group while there was no survival difference between the UFT and the control group in the patients with T1 disease. The 5-year survival rate of patients with T2 disease was 85 percent (95 percent confidence interval: 79 to 91 percent) in the UFT group and 74 percent (95 percent confidence interval: 68 to 81 percent) in the control group ( FIG. 1B ). The overall survival between the two groups was statistically significantly different (P=0.005 by the logrank test). The 5-year survival rate of patients with T1 disease was 89 percent in the UFT and 90 percent in the control groups ( FIG. 1C ). The 5-year survival rate of patients with tumor size of <=2 cm, 2 to 3 and >3 was 89, 89 and 85 percent in the UFT group and 91, 86 and 74 percent in the control group, respectively. Pattern of Failure and Cancer-Free Survival As shown in Table 3, either recurrence or a second primary cancer as the first treatment failure after operation was documented in 23 percent of the patients in the UFT group and in 26 percent of those in the control group. Among 716 patients with T1 disease, either recurrence or a secondary primary cancer as the first treatment failure was observed in 69 patients (19 percent) in the UFT group and 76 patients (22 percent) in the control group while it was observed in 42 patients (33 percent) in the UFT group and 53 patients (40 percent) in the control group among 263 patients with T2 disease. The cancer-free survival between the UFT and the control group was not statistically significantly different based on a Kaplan-Meier analysis (P=0.25 by the stratified logrank test). The survival of patients after the diagnosis of either recurrence or second primary cancer did not differ significantly between the groups (P=0.14 by the logrank test): the 1- and 2-year survival rates after diagnosis were 65 percent and 50 percent in the UFT group and 65 percent and 42 percent in the control group, respectively. The Japanese Association for Chest Surgery and Japan Lung Cancer Society recently reported on the long-term survival rate of 7,408 lung cancer patients who underwent a surgical resection in 1994 when the present trial started. (Shirakusa, T., et al., 2002). The main histologies were adenocarcinoma (56 percent) and squamous cell carcinoma (33 percent). The 5-year survival rate of patients with pathological T1N0M0 and T2N0M0 was 79 percent and 60 percent, respectively. In the present study focusing on the histology of adenocarcinoma, the 5-year survival rate in the control group was 90 percent in patients with T1N0M0 and 74 percent in those with T2N0M0. Although the latter figures could not be directly compared with the former due to different histologic patterns and the times that the data were collected, the 5-year survival of the control patients in the present study is thought to be at least one of the best results reported. (Breathnach, O. S., et al., 2001; Myrdal, G., et al., 2003). Those findings may indicate the quality of the operation and the accuracy of surgical staging in the study group. It was thus shown that postoperative adjuvant chemotherapy with UFT has a beneficial effect on the survival of resected patients with stage I adenocarcinoma, although little or no benefit was observed in the patients with T1N0 disease. Recently, the number of patients with small sized adenocarcinoma has increased due to the increased use of CT. In fact, 412 (42 percent) of total 979 patients had adenocarcinoma measuring less than 2 cm in size. Such small sized adenocarcinomas often include bronchioloalveolar carcinoma which has little chance of recurrence after operation. (Noguchi, M., et al., 1995). As a result, prognosis of adenocarcinoma measuring less than 2 cm in size is very good (Noguchi, M., et al., 1995; Kodama, K., et al., 2001): the 5-year survival rate in the present study was 91 percent. Therefore, the patients of this group should be excluded from any postoperative adjuvant trials in the future if a poor prognosis subgroup cannot be identified. In contrast, patients with a tumor size ranging from 2 to 3 cm tended to show an improved survival by UFT treatment while those with a tumor size of over 3 cm who received UFT treatment had a definitive beneficial effect on survival. These findings indicate that the effect of UFT may be related with some biological malignant factors. Tanaka et al. (Tanaka, F., et al., 2001) reported that patients with non-small cell lung cancer with a high apoptotic index and no p53 aberrant expression who underwent postoperative chemotherapy with UFT demonstrated good prognosis in a retrospective study. Patient compliance in postoperative adjuvant trials is always a problem. In the trials using cisplatin-based chemotherapy, which was planned to be administered in three or four cycles after operation, only 50 to 70 percent of the planned treatment has been reported to be successfully performed. (Feld, R., et al., 1993; Ohta, M., et al., 1993; Ichinose, Y., et al., 2001; Scagliotti, G. V., et al., 2003) In the present trial, UFT was planned to be given daily for two years. However, only 61 percent of candidate patients completed the two-year treatment in spite of mild toxicities. In fact, the main reasons for discontinuing UFT administration were adverse reactions and patient refusal. Those findings indicate that the compliance of adjuvant chemotherapy trials may not be related to the severity of adverse events induced by chemotherapy. However, if the effect of adjuvant chemotherapy is proven, it is clear that a treatment with mild toxicity has better compliance than a treatment with severe toxicity. The main difference between cisplatin-based adjuvant chemotherapy trials and UFT adjuvant trials is the duration of the treatment. The former trials have three or four cycles (9 to 16 weeks) (Feld, R., et al., 1993; Ohta, M., et al., 1993; Ichinose, Y., et al., 2001; Scagliotti, G. V., et al., 2003) whereas UFT has been administered daily for 1 or 2 years in the latter trials. (Wada, H., et al., 1996; Tanaka, F., et al., 2001; Tada, H., et al., 2002; Endo, C., et al., 2003; Imaizumi, M., et al., 2003). 5-fluorouracil is well known to not be a dose-dependent agent, but to be a time-dependent agent. Therefore, the daily administration of UFT is an effective method for maintaining the 5-fluorouracil concentration in the blood. In addition, UFT and its metabolites have both been recently reported to have an inhibitory effect on tumor angiogenesis. (Yonekura, K., et al., 1999). If this effect is truly present in the human body, then the daily and long-term administration of UFT may be an even more ideal administration method. So far, six randomized trials (Wada, H., et al., 1996; Tanaka, F., et al., 2001; Tada, H., et al., 2002; Endo, C., et al., 2003; Imaizumi, M., et al., 2003), including the present trial, comparing surgery alone with postoperative adjuvant treatment with UFT, have been conducted. Among them, three trials demonstrated a survival benefit of UFT. (Wada, H., et al., 1996; Tada, H., et al., 2002). In addition, the results of a meta-analysis of those six trials demonstrated that adjuvant chemotherapy with UFT improved the overall survival (hazard ratio=0.77, 95 percent confidence interval=0.63 to 0.94 percent; P=0.01). (Hamada, C., et al., 2003). Whether the patients with stage II or III have a survival benefit from UFT treatment or whether the 1-year treatment is equivalent to 2-year treatment remains unclear. Nonetheless, based on the results of the present study, patients with completely resected stage I disease, especially T2N0 adenocarcinoma, should benefit from long term, postoperative adjuvant chemotherapy with UFT for a period of 2 or more years. The present invention is further illustrated by the following example which is not intended to be limiting. EXAMPLE 1 Patients who had undergone a complete resection of a pathologically documented stage I (T1–2, N0, M0) (Mountain, C. F., et al., 1986) adenocarcinoma were eligible for treatment in accordance with the present invention. Visceral pleural involvement was classified according to the rules of the Japan Lung Cancer Society (Japan Lung Cancer Society, 1987) and either a tumor that was larger than 3 cm or a tumor with any size that was exposed on the visceral pleural surface was classified as T2 tumor. Patients were assigned to either a treatment or control group. In the treatment group (n=491), UFT (tegafur 250 mg per square meter of body-surface area) in the form of a 100-mg capsule (100 mg tegafur and 224 mg uracil) was given orally in two separate doses, before meals, daily for two years, starting four weeks after operation. The dose was rounded up or down to the nearest 100 mg. Most patients received UFT as two capsules (tegafur 200 mg and uracil 448 mg) bis in die. On each visit at the outpatient clinic, the physician in charge asked all patients whether they regularly took UFT capsules as ordered. Patients assigned to the control group (n=488) were observed with no further treatment after operation. The toxicity resulting from UFT administration was graded according to the criteria of the Japan Society of Clinical Oncology, which consist of the World Health Organization criteria with minor modifications. (World Health Organization, 1979). If a grade 2 adverse reaction occurred, the dose of UFT was reduced to 200 mg per square meter. If a grade 3 or greater adverse reaction, a leukocyte count of less than 3,000 per cubic millimeter, a platelet count of less than 70,000 per cubic millimeter, a hemoglobin level of less than 9.5 g per deciliter or aspartate aminotransferase and alanine aminotransferase levels that were more than three times the upper limit of the normal range occurred, then the administration of UFT was suspended. A follow-up examination was performed every three months for the two years after the patients operation and every six months thereafter. The examination included a physical examination, a complete blood count, blood chemistry work-up, serum tumor marker screening, and chest radiography. A computed tomographic (CT) scan of the thorax and brain, and either a CT scan or an echogram of the upper abdomen were performed every six months for the first two years after the patient's operation and at least twice during the subsequent three years. Any newly appearing lesion suspected of either being recurrence or a second primary cancer was investigated by a biopsy whenever possible. A final diagnosis of such lesions as either recurrence or a second primary cancer was made by the physician in charge. The primary end point was the overall survival, while the secondary end points were cancer-free survival and safety assessment. The subjects included in the analysis of overall survival and cancer-free survival were all eligible patients. The subjects in the safety assessment consisted of patients who were given UFT. The sample size was calculated by the method of Schoenfeld and Richter (Schoenfeld, D. A., et al., 1982) under the following conditions: a 5-year survival rate of 70 percent for the non-treatment control group, a hazard ratio of 0.667 for death in the UFT group, the 2-year accrual period, the 5-year follow-up, a significance level for a one-sided test of 0.05, and a statistical power of 80 percent. Since the above calculations resulted in a sample size of 518 patients, the sample size was determined to be 600, with an allowance of about 15 percent for ineligible cases or cases which were lost to follow-up. In May 1995, the sample size was expanded to 984 patients, because it became clear that the 5-year survival rate for those in the control group was better than expected. The newly adopted survival rate was 83 percent, and the accrual period was extended to 3 years. Committee for Efficacy and Safety provided independent monitoring of the study. Haybittle-Peto horizontal boundaries, (Haybittle, J. L., 1971) with a criterion of P<0.001, were used in the interim analyses conducted to determine whether the study should be terminated early. The overall survival was defined as the time from operation until death from any cause, and the cancer-free survival was defined as the time from operation until the appearance of either the first recurrence of cancer, a second cancer or death from any cause. Survival was estimated by the Kaplan-Meier method, and any differences in survival were computed using the stratified log-rank test. Multivariable analyses using the Cox proportional hazard model were used to estimate the simultaneous effects of prognostic factors on survival. (Cox, D. R., 1972). The interactions with prognostic factors were examined with the Cox proportional hazard model. The SAS statistical software package was used for all calculations. The data were considered to be statistically significant when the P value was 0.05 or less. All statistical tests were two-sided. The median follow-up for surviving patients was 73 months. The overall survival between the two groups showed a statistically significant difference in favor of the UFT group based on a Kaplan-Meier analysis (P=0.04) by the stratified logrank test. Grade 3 toxicity based on UFT administration was observed in only 10 (2 percent) of 482 patients. The publications and other materials cited herein to illuminate the background of the invention and to provide additional details respecting the practice of the invention are incorporated herein by reference to the same extent as if they were individually indicated to be incorporated by reference. While the invention has been disclosed by reference to the details of preferred embodiments of the invention, it is to be understood that the disclosure is intended in an illustrative rather than a limiting sense, as it is contemplated that modifications will readily occur to those skilled in the art, within the spirit of the invention. REFERENCES 1. Fujii, S., et al., “Effect of coadministration of uracil or cytosine on the anti-tumor activity of clinical doses of 1-(2-tertahydrofuryl)-5-fluorouracil and level of 5-fluorouracil in rodents,” Gann 70:209–14, 1979. 2. Ikenaka, K., et al., “Effect of uracil on metabolism of 5-fluorouracil in vitro,” Gann 70:353–9, 1979. 3. Ho, D. H., et al., “Comparison of 5-fluorouracil pharmacokinetics in patients receiving continuous 5-flourouracil infusion and oral uracil plus N1-(2′-Tertahydrofuryl)-5-flourouracil,” Clin. Cancer Res. 4:2085–8, 1998. 4. Shimizu, E., et al., “A phase II study of UFT in non-small cell lug cancer,” Jpn. J. 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Okimoto, N., et al., “A randomized controlled postoperative adjuvant chemotherapy trial of CDDP+VDS+UFT and UFT alone in comparison with operation only for non-small cell lung carcinomas (Second Study),” Jpn. J. Lung Cancer 36:863–71, 1996. 15. Mountain, C. F., “A new international staging system for lung cancer,” Chest. 89(4):225s–33s, 1986. 16. WHO Handbook for Reporting Results for Cancer Treatment , Geneva: World Health Organization Offset Publication No. 48, 14–21, 1979. 17. Shirakusa, T., et al., “Lung cancer in Japan: Analysis of lung cancer registry for resected cases in 1994,” Jpn. J. Lung Cancer 42:555–66, 2002. 18. Breathnach, O. S., et al., “Bronchioloalveolar carcinoma of the lung: recurrences and survival in patients with stage I disease,” J. Thorac. Cardiovasc. Surg. 121:42–7, 2001. 19. Myrdal, G., et al., Survival in primary lung cancer potentially cured by operation: influence of tumor stage and clinical characteristics,” Ann. Thorac. Surg. 75:356–63, 2003. 20. 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Ichinose, Y., et al., “A Randomized Phase III Trial of Postoperative Adjuvant Chemotherapy in Patients with Completely Resected Stage IIIa-N2 Non-Small Cell Lung Cancer: Japan Clinical Oncology Group (JCOG9304) Trial,” Prog. Proc. Am. Soc. Clin. Oncol. 20:311a (Abstract), 2001. 26. Scagliotti, G. V., et al., “Randomized study of adjuvant chemotherapy for completely resected stage I, II, or IIIA non-small-cell lung cancer,” J. Natl. Cancer Inst. 95:1453–1461, 2003. 27. Tanaka, F., et al., “Postoperative oral administration of UFT for completely resected pathologic stage I non-small cell lung cancer: the West Japan Study Group for Lung Cancer Surgery (WJSG), the 4 th Study,” Prog. Proc. Eur. Cancer Conference 37:S29 (Abstract), 2001. 28. Tada, H., et al., “Randomized Study of Adjuvant Chemotherapy for Completely Resected Non-Small Cell Lung Cancer,” Prog. Proc. Am. Soc. Clin. Oncol. 21:313a (Abstract), 2002. 29. Endo, C., et al., “A randomized trial of postoperative UFT therapy in p stage I, II non-small cell lung cancer: North-East Japan Study Group for Lung Cancer Surgery,” Lung Cancer 40:181–6, 2003. 30. Imaizumi, M., et al., “A randomized trial of postoperative adjuvant chemotherapy for p-stage I non-small cell lung cancer (4 th cooperative study),” Prog. Proc. World Conference 41:S54 (Abstract), 2003. 31. Yonekura, K., et al., UFT and its metabolites inhibit the angiogenesis induced by murine renal cell carcinoma, as determined by a dorsal air sac assay in mice,” Clin. Cancer Res. 5:2185–91, 1999. 32. Hamada, C., et al., “Efficacy or oral UFT for adjuvant chemotherapy after complete resection of non-small cell lung cancer: Meta-analysis of-six randomized trials in 2003 patients,” Prog. Proc. Euro. Cancer Conference 39:S231 (Abstract), 2003. 33. Mountain, C. F., “Revisions in the International System for Staging Lung Cancer,” Chest 111 (6):1710–17, 1997. 34. Mountain, C. F., et al., “Regional Lymph Node Classification for Lung Cancer Staging,” Chest 111 (6):1718–1723, 1997. 35. The Japan Lung Cancer Society, General Rule for Clinical and Pathological Record of Lung Cancer, 3rd Ed. Tokyo: Kanehara, 1987. 36. Schoenfeld, D. A., et al., “Nomograms for calculating the number of patients needed for a clinical trial with survival as an endpoint,” Biometrics 38:163–170, 1982. 37. Haybittle, J. L., “Repeated assessment of results in clinical trials of cancer treatment,” Br. J. Radiol. 44:793–797, 1971. 38. Cox, D. R., “Regression models and life-tables,” J. Roy Statist. Soc. Ser. B. 34:187–220, 1972. 39. Kato, H., et al., “A randomized phase III trial of adjuvant chemotherapy with UFT for completely resected pathological stage I (T1N0M0, T2N0M0) adenocarcinoma of the lung,” Proc. Am. Soc. Clin. Oncol. 22:621, 2003 (Abstract 2498). TABLE 1 Patient Characteristics Control UFT Characteristics (N = 488) (N0 = 491) Age (yr) Mean (range) 62 (45–75) 62 (45–75) <65 275 274 >65 213 217 Female sex (no. of patients) 249 253 ECOG performance status (no. of patients) 0 369 376 1 113 105 2 6 10 Pathological T status (no. of patients) 1 354 362 2 134 129 Pleural invasion (no. of patients)* 0 346 340 1 114 120 2 28 29 Unknown 0 2 Tumor size (no. of patients) <=2 cm 204 208 2 to 3 cm 170 174 >3 cm 114 109 Location of the tumor (no. of patients) Rt. upper lobe 189 182 Rt. middle lobe 34 41 Rt. lower lobe 87 102 Rt. Lobes 2 2 Lt. upper lobe 114 107 Lt. lower lobe 60 54 Lt. lobes 2 3 Operation modality (no. of patients) Lobectomy 487 490 Pneumonectomy 1 1 *0 = a tumor with no pleural involvement or a tumor that reaches the visceral pleura but does not extend beyond the elastic layer; 1 = a tumor that extends beyond the elastic layer of the visceral pleura but is not exposed on the pleural surface; 2 = a tumor that is exposed on the pleural surface but does not involve the parietal pleura. TABLE 2 UFT-Related Adverse Reaction Toxicity (n = 482), Grade percent of patients 1 2 3 4 Leukopenia 2 1 0 0 Thrombocytopenia <1 0 0 0 Hemoglobin <1 <1 0 0 Bilirubin 1 <1 0 0 GOT 6 2 <1 0 GPT 6 2 0 0 ALP 2 <1 0 0 Anorexia 9 8 1 — Nausea/Vomiting 10 3 1 — Diarrhea 2 1 <1 0 Alopecia <1 0 0 — TABLE 3 Pattern of Treatment Failure Control UFT Pattern (n = 488) (n = 491) Intrathoracic only Local recurrence 8 17 Pulmonary metastases 38 36 Local recurrence plus Pulmonary metastases 12 3 Second Cancer 11 11 Extrathoracic only Recurrence 33 23 Second cancer 18 14 Intrathoracic plus Extrathoracic recurrence 9 7 Total No. (percent of all patients) 129 111 (26.4) (22.6)

The present invention provides an improved method for treating lung cancer, preferably non-small cell lung cancer, by orally administering UFT to postoperative lung cancer patients.

FIELD OF THE INVENTION The present invention relates to a process for recovering fluoride from fluoride-containing wastewater. More specifically, it relates to a process for recovering fluoride as calcium fluoride from fluoride-containing wastewater, for example, produced in the manufacture of crystalline silicon solar cell. The recovered calcium fluoride has high purity and large particle sizes. Therefore, it may be used as a raw material for producing hydrogen fluoride or for other applications. BACKGROUND OF THE INVENTION Many industrial operations generate fluoride-containing wastewater. For example, even though the process of generating green photovoltaic (PV) electricity from solar cells creates no pollution or waste, the manufacture of solar cells can produce serious pollutions. The fluoride-containing wastewater accounts for a very large percentage of the total hazardous wastes generated in crystalline silicon solar cell manufacturing. The fluoride-containing wastewater must be treated to comply with the industrial effluent discharge standards. For clarity, the following description may use the production of crystalline silicon solar cells as an example to illustrate the invention. However, one skilled in the art would appreciate that embodiments of the invention can also be used with fluoride-containing wastewater generated in other operations. Crystalline silicon solar cells are made from thin silicon wafers or slices, and the manufacture process consists of six main steps, viz, cleaning, texturization, emitter formation, antireflection film deposition, metallization, and cell testing, as shown in FIG. 1 . These six steps are briefly described as follows: Cleaning Raw silicon wafers are etched in hydrochloric acid (HCl) baths to remove organics, silicon carbide and other impurities produced from the wafering processes. The HCl etching process leaves a waste hydrochloric acid pickling liquor. Texturization Texturization is a process to etch the wafers to remove any residual saw damages and to form surface textures to improve crystalline silicon solar cell efficiencies by lowering surface reflection and increasing light absorption. Anisotropic alkaline texturing is mainly used for mono-crystalline silicon wafers, while isotropic acidic texturing is mainly used for multi-crystalline silicon wafers. In alkaline texturing, silicon wafer is typically textured in a mixture of sodium hydroxide (NaOH) and isopropanol (IPA) at 80° C. After texturing, the wafers are rinsed in deionized (DI) water, cleaned in hydrochloric acid (HCl) solution, rinsed in DI water, cleaned in hydrofluoric acid (HF) solution, rinsed in DI water, and then dried in hot air. In this process, HCl is used to remove metal impurities from the wafer surface, while HF is used to strip native silicon dioxide and to form a hydrophobic surface. This process produces a waste NaOH texturing liquor, a waste HCl acid pickling liquor, and a waste HF acid pickling liquor, and. In acidic texturing, silicon wafer is generally textured in a mixture of hydrofluoric acid (HF) and nitric acid (HNO 3 ), wherein HNO 3 , an oxidizing agent, is used to form a silicon oxide on the surface of the wafer. After texturing, the wafers are rinsed in deionized (DI) water, cleaned in alkaline (NaOH) solution, rinsed in DI water, cleaned in hydrofluoric acid (HF) solution, rinsed in DI water, and then dried in hot air. This process leaves a waste HNO 3 /HF acid pickling liquor and a waste NaOH rinsing liquor. Emitter Formation In a p-type silicon solar cell, the emitter layer is formed by thermal diffusion of phosphoric oxychloride (POCl 3 ) to form the pn junction. This thermal diffusion is typically performed in a diffusion oven at a temperature of 800-900° C. 4POCl 3 +3O 2 →2P 2 O 5 +6Cl 2 ↓ In this process, P 2 O 5 forms a layer of phosphor silicate glass (PSG) on the wafer, from which P atoms will diffuses into the upper part of the wafer. Meanwhile, after the diffusion step, the edges of the wafer also contain P atoms, and the edges must be isolated in order to prevent short-circuiting. PSG removal and edge isolation are completed by etching with a solution of HF and HNO 3 . SiO 2 +6HF→H 2 SiF 6 +2H 2 O After PSG removal and edge isolation, the wafers are rinsed in deionized (DI) water, cleaned in sodium hydroxide (NaOH) solution, rinsed in DI water, cleaned in HF/HNO 3 solution, rinsed in DI water, and then dried in hot air. This process leaves a waste HNO 3 /HF acid pickling liquor, as well as a waste NaOH rinsing liquor. Antireflection Film Deposition The antireflection coating is carried out by plasma enhanced chemical vapor deposition (PECVD) of Si 3 N 4 by the reaction between silane (SiH 4 ) and ammonia (NH 3 ). The exhaust of the PECVD reactor is connected with a silane burner, where excess silane is converted to SiO 2 and H 2 O. Metallization A conductive pattern is screen printed onto the front wafer surface with a silver-rich paste. The front conductive pattern forms the busbar that will conduct the photo-electrons. A back print busbar is screen printed onto the cell's rear surface using a silver-rich paste. The rear busbar provides a means for soldering and interconnecting adjacent cells to make a solar panel. Finally, the entire rear surface is screen printed with aluminium to create a back field surface (BSF) and to enhance gettering. After the paste has been deposited, the wafers are fired in a belt oven. Emissions that occur during the firing process are evaporated solvents and burnt organic compounds. Cell Testing The cells are individually tested for electrical performance under Standard Test Conditions (STC) of 1 kW/m 2 (1 Standard Sun) irradiance and 25° C. Subsequently, the cells are graded into specific power bands for future cell stringing to make panels. This process does not produce wastes. Among the processes described above, texturization, PSG removal, and edge isolation are typically performed in batch processes using inline equipment. Wafers are held in cassettes that allow chemicals to wet the entire surfaces of the wafers. These cassettes are moved automatically on rolls from one tank to another filled with chemicals for etching, cleaning, or rinsing. The production process and the process parameters of different silicon solar cell manufacturers may vary. However, the produced wastewaters are essentially the same, which all consist of acid wastewater, alkaline wastewater, and rinsing wastewater. The main characteristics of these wastewaters are high acidity, high basicity, or high fluoride concentrations. The environmentally toxic pollutant in the wastewater is fluoride because of its demonstrated link to human fluorosis. Therefore, the fluoride-containing wastewater must be treated to meet the wastewater discharge standards of the local countries. For example, according to Chinese regulations for industrial wastewater, the permissible standard is F − ≦20 mg/l. Until now, most crystalline silicon solar cell manufacturers use the precipitation-flocculation process for the removal of fluoride due to its relative simplicity and low costs. The process is shown in FIG. 2 . As illustrated in FIG. 2 , the acidic wastewater and alkaline wastewater are discharged, respectively. Then, the traditional precipitation-flocculation method is used to remove fluoride. This process consists of neutralization, chemical precipitation, flocculent precipitation, and filtration. The filtered calcium fluoride sludge is discarded. In the precipitation-flocculation process, either calcium chloride (CaCl 2 ) or calcium hydroxide (Ca(OH) 2 ) is used to induce the precipitation of fluoride as calcium fluoride (CaF 2 ), according to the following chemical reactions: Ca(OH) 2 +2HF→CaF 2 +2H 2 O CaCl 2 +2HF→CaF 2 +2HCl When Ca(OH) 2 is used, it should be significantly overdosed to work effectively due to its low solubility. As a result, a larger amount of sludge is generated, and the unreacted sparsely soluble Ca(OH) 2 particles become the major constituent in the precipitated sludge. Therefore, CaCl 2 is the preferred choice due to its higher solubility, and a smaller amount of sludge is generated when CaCl 2 is used. However, when CaCl 2 is used, the very fine CaF 2 particles do not settle readily. Therefore, it becomes necessary to use flocculants, such as polyaluminum chloride (PAC) or polyacrylamide (PAM), to help bring down (flocculate) the fine CaF 2 particles for easier separation from the wastewater. In practice, many treatment plants use a combination of Ca(OH) 2 and CaCl 2 to enhance the settlement of CaF 2 precipitates, though a larger amount of sludge is produced, as compared to using only CaCl 2 . Similar precipitation-flocculation processes are described in U.S. Pat. No. 7,182,873 B2, Chinese Patent No. CN 101973662, Worf et al., “waste water treatment for crystalline silicon solar cell production,” Photovoltaics International, 11th Edition, and Drouiche et al., “photovoltaic solar cells industry wastewater treatment,” Desalination and Water Treatment, 2013, pp 1-9. In these processes, the fluoride-containing wastewater may be introduced into a neutralization tank to neutralize it by adding NaOH or HCl. The neutralized wastewater is then delivered to a reaction tank to add CaCl 2 and/or Ca(OH) 2 . Then, the wastewater and the generated calcium fluoride are delivered to a settling tank, therein agglutinators or flocculants are added to separate the calcium fluoride from wastewater. The sludge generated in the precipitation-flocculation processes has a moisture content of as much as 50-60% and may contain a wide variety of impurities, such as Al 3+ , Cl − , SO 4 2− , unreacted Ca(OH) 2 , etc. Therefore, the sludge cannot be used in other industrial applications. Nowadays, the CaF 2 sludge is typically transported to controlled hazardous waste disposal sites. Due to huge volumes, the disposal costs are high. Furthermore, the sludge has a potential to cause secondary pollution. For example, if the environment has acid rain or acidic soil, the fluoride ion in the sludge will enter into the groundwater and pollute the soil and water. On the other hand, calcium fluoride is widely used as a feedstock in the manufacture of HF and as a fluxing agent in the steel industry, as well as in the production of glass, enamel, welding rod coatings, etc. Currently, almost all fluorine or fluoride used in various industries is processed from fluorspar ores. However, fluorspar is becoming a more and more precious non-renewable resource. Therefore, there is an urgent need to develop technologies for recovering fluoride from industrial fluoride-containing wastewater. SUMMARY OF THE INVENTION An object of the present invention is to provide a process for recovering fluoride as calcium fluoride from fluoride-containing wastewater, such as the fluoride-containing wastewater produced in the manufacture of crystalline silicon solar cell. More particularly, the present invention provides economically and ecologically beneficial treatment processes that can be used to recover, with high efficiencies (e.g., over 70%, over 80%, or over 90%), fluoride from a fluoride-containing wastewater. The fluoride-containing wastewater includes that produced in the manufacture of crystalline silicon solar cell or in other industrial operations. Embodiments of the invention may greatly reduce the amount of calcium fluoride sludge waste (for example, to zero) by producing calcium fluoride of higher purity (e.g., over 90%, over 95%, or over 98% pure) that can be used in other operations. A process in accordance with embodiments of this invention may comprise four sections A, B, C and D. Section A may be referred to as “separate collection.” Section B may be referred to as “preparation of calcium reagent solution.” Section C may be referred to as “primary treatment.” Section D may be referred to as “secondary treatment.” Section A comprises classifying and collecting the wastewater from a source, e.g., from the manufacture of crystalline silicon solar cell, into a few types (e.g., 3 types) of streams, according to the concentrations of fluoride and the pH values. For example, 3 types of streams may be, respectively, alkaline wastewater, acidic wastewater with a high fluoride concentration, and acidic wastewater with a low fluoride concentration. Whether a particular wastewater is considered as having a high or low fluoride concentration may be based on a criterion (a threshold), which might be judicially selected based on the source of the wastewater. For example, for wastewater from crystalline silicon solar cell production, this criterion (threshold), for example, may be 1 g/L, 5 g/L, or 10 g/L. One skilled in the art would know that this criterion may be different for different sources of wastewater, and consideration may be given for efficient handling and treatment of the wastewater. Similarly, what is classified as high pH or low pH may be judicially determined based on another threshold. Typically, wastewater generated from the production of crystalline silicon solar cells may have a relatively high pH (e.g., >9.0, >10.0, or >11.0) or a relatively low pH (e.g., <5.0, <3.0, or <2.0). If desired the pH value of wastewater may be adjusted using a base or acid to a desired level. In Section B, the acidic wastewater with a low fluoride concentration from Section A may be reacted with calcium carbonate (CaCO 3 ), calcium oxide (CaO) or calcium hydroxide (Ca(OH) 2 ) to prepare an aqueous solution containing Ca 2+ ion. This calcium-containing solution may be used to precipitate fluoride in the following processing steps. In Section C, the aqueous solution containing Ca 2+ ion from Section B and the acidic wastewater with a high fluoride concentration from Section A are sent to the first reactor, wherein calcium fluoride is precipitated. Ca 2+ 2F − →CaF 2 ↑ The calcium fluoride precipitates are separated from the wastewater. The separation may use any equipment known in the art, such as a solid-liquid separator. The effluent from the separator would have a lower fluoride concentration and may be mixed with the rest of the acidic wastewater with a low fluoride concentration in a mixer tank. The reactor used in Section C may be a stirred tank reactor (STR) or any other suitable reactors. In Section D, the mixed acidic wastewater from Section C and the prepared aqueous solution containing Ca 2+ ion from Section B are thoroughly mixed within the second reactor containing a suitable seed material such that the precipitated calcium fluoride can crystallize. In this treatment, the acidity of the mixed acidic wastewater may be adjusted by combination with the alkaline wastewater. The reactor used in this section may be a fluidized bed reactor (FBR) or any other suitable reactors. The CaF 2 particles produced from the first reactor in Section C may be used as seed materials. Alternatively, calcium fluoride particles or crystals from other sources may be used. The FBR may use an upflow of a liquid or diffused air to ensure that the mixture in the reactor remains homogeneous. The FBR unit may also comprise a clarifier, an aeration blower, and a delivery system for various liquid streams. The above describes a preferred embodiment of the invention that includes four sections, in which sections C and D both produce calcium fluoride particles. One skilled in the art would know that Section C may precipitate most of the fluoride in the wastewater, and the effluent from Section C may contain relatively low concentrations of fluoride. In some situations, it may not be necessary to further remove fluoride from the “low fluoride-containing” water. In those cases, one may skip Section D, even though inclusion of Section D would produce more calcium fluoride and cleaner (containing less fluoride) wastewater. Embodiments of the invention will now be illustrated with a specific process with more details. One skilled in the art would know that this particular example is for illustration only and that variations and modifications of various details are possible without departing from the scope of the invention. For example, a process for processing fluoride-containing wastewaters from a factory (e.g., the manufacture of crystalline silicon solar cells) in accordance with one embodiment of the invention may include the following steps: Step 1: collecting the fluoride-containing wastewaters into three pools: an acidic high-fluoride wastewater that contains fluoride at a concentration higher than a first threshold concentration, an acidic low-fluoride wastewater that contains fluoride at a concentration lower than a second threshold concentration, and an alkaline wastewater. As used herein, the term “acidic” refers to a pH value lower than 7, preferably lower than 5, such as 3, or 2, or lower, and the term “alkaline” refers to a pH value greater than 7, preferably greater than 9, such as pH 10 or 11 or higher. As used herein, the term “high-fluoride” refers to a fluoride concentration higher than a threshold value, wherein the threshold value may be 1000 mg/L, preferably 2,000 mg/L, and more preferably 5,000 mg/L. As used herein, the term “low-fluoride” refers to a fluoride concentration lower than a threshold value, wherein the threshold value may be 2,000 mg/L, preferably 1,000 mg/L, more preferably 500 mg/L, and most preferably 100 mg/L. Step 2: adding a calcium compound to a portion of the acidic low-fluoride wastewater to produce a “calcium-containing solution” that contains calcium at a concentration of about 1-500 g/L, preferably about 10-300 g/L, and more preferably about 80-120 g/L. In accordance with embodiments of the invention, the calcium compound may be calcium carbonate or any suitable calcium salt (e.g., Ca(OH) 2 , CaCl 2 , or CaO). The calcium compound (e.g., calcium carbonate) may be added as powders or as a suspension or solution in water at a suitable concentration (e.g., 5-50%). Step 3: reacting a portion of the calcium-containing solution with the acidic high-fluoride wastewater at a suitable calcium-to-fluoride molar ratio, e.g., at a ratio of from about 0.5:1 to about 1.5:1, to produce a mixture comprising calcium fluoride particles suspended in a mother liquor. The term “molar ratio” as used herein has its regular meaning, i.e., a ratio of the moles of calcium and the moles of fluoride. The term “mother liquor” is used as its regular meaning, i.e., the liquid after precipitates are removed. Step 4: processing the mixture to separately collect the calcium fluoride particles and the mother liquor. The processing involves separation of precipitates (CaF 2 particles) from the mother liquor. The mother liquor would have a lower concentration of fluoride ions. The separation can be performed with any method and equipment known in the art, such as gravitational settlement, filtration, or centrifugation. A suitable equipment may be a sedimentation centrifuge, a filtration centrifuge, or a pressure or suction filtration machine. Step 5: diluting the mother liquor with a diluent to produce a mixed solution having a suitable pH value (e.g., about pH 2-6) and containing a suitable fluoride concentration (e.g., about 80-500 mg/L). The diluent used here may be any liquid (e.g., water) or solution with low or no fluoride ions. A convenient diluent for this, for example, may be the acidic low-fluoride wastewater obtained in Step 1. Using the acidic low-fluoride wastewater from Step 1 is advantageous because it would minimize the total volume of wastewater that needs to be disposed of. Step 6: introducing the mixed solution from Step 5, the calcium-containing solution from Step 2, and the alkaline wastewater from Step 1 into a fluidized-bed reactor, which contains a calcium fluoride crystallization seed material, to form calcium fluoride crystals, wherein the alkaline wastewater is used to adjust a pH of the fluid in the fluidized-bed reactor (FBR) to a value between 5 and 8. The use of the alkaline wastewater to adjust pH has the advantage of minimizing the total volume of wastewater that needs to be disposed of, even though other source of alkaline solution may be optionally used here. In this step, a molar ratio between the calcium in the calcium-containing solution from Step 2 and the fluoride in the mixed solution from Step 5 is preferably controlled at a suitable ratio, such as 0.5:1 to 2:1. In this Step 6, formation of calcium fluoride particles or crystals may be optimized at a suitable fluoride concentration and a suitable fluid flow rate in the FBR. The concentration of fluoride may be adjusted using water at the inlet. Alternatively, the effluent from the fluidized bed reactor in this Step 6 may be used. The effluent would have relatively lower concentration of fluoride ions and, therefore, can be used for this purpose. Using the effluent for this purpose has the advantage of minimizing the total volume of wastewater that needs to be disposed of. A suitable fluoride concentration, which may be determined by any suitable methods (e.g., by trial runs), for example may be at around 150 mg/L or lower. One skilled in the art would also know that a suitable flow rate in the fluidized bed reactor may be determined by any suitable methods (e.g., by trial runs). As an example, the flow rate in FBR may be controlled at about 100-20,000 m/hr. A process in accordance with embodiments of this invention may have one or more of the following advantages: 1. markedly reduced sludge volumes and therefore reduced disposal costs; 2. recovery of the valuable resource CaF 2 ; 3. increased process stability and reliability; and 4. better effluent quality. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 shows a flowchart illustrating a process flow diagram for the manufacturing of crystalline silicon solar cells. FIG. 2 shows a flowchart illustrating a prior art precipitation-flocculation process for the treatment of fluoride-containing wastewater. FIG. 3 shows a schematic illustration of a system for recovering fluoride from wastewater produced in the manufacture of crystalline silicon solar cells according to one embodiment of the present invention. FIG. 4 shows a schematic illustration of a fluidized bed reactor in accordance with one embodiment of the present invention. DETAILED DESCRIPTION OF THE INVENTION Embodiments of the invention relate to methods (processes) for treating fluoride-containing wastewater. These processes not only remove the harmful fluoride from the wastewater, but also produce useful calcium fluoride for other applications. A process in accordance with embodiments of this invention may comprise four sections (or sub-processes) A, B, C and D. Section A may be referred to as “separate collection”, Section B may be referred to as “preparation of calcium reagent solution”, Section C may be referred to as “primary treatment,” and Section D may be referred to as “secondary treatment.” Referring to FIG. 3 , in Section A, a wastewater (for example, a wastewater produced from the manufacture of crystalline silicon solar cells) may be classified into three types: an acidic wastewater with a low fluoride concentration (S1), an acidic wastewater with a high fluoride concentration (S2), and an alkaline wastewater (S3). The three types of wastewater may be stored, respectively, in a dilute acidic wastewater tank A1, a concentrated acidic wastewater tank A2, and an alkaline wastewater tank A3. In accordance with embodiments of the invention, an acidic wastewater will have a pH value lower than 7. Because the acidic wastewater is typically from HCl wash or HF etch, its pH value is typically much lower than 7, such as 3 or lower. In accordance with embodiments of the invention, an alkaline wastewater will have a pH value higher than 7. Typically, a wastewater with a pH value over 9 may be collected as an alkaline wastewater. In accordance with embodiments of the invention, an acidic wastewater with a fluoride concentration over 1,000 mg/l, preferably over 2,000 mg/l, more preferably over 5000 mg/l, may be collected as an acidic wastewater with a high fluoride concentration, while an acidic wastewater with a fluoride concentration below 1,000 mg/l, preferably below 500 mg/l, more preferably below 100 mg/l, may be collected as an acidic wastewater with a low fluoride concentration. For instance, a wastewater from the cleaning step in the production of crystalline silicon solar cells contains only HCl. Therefore, it may be collected as an acidic wastewater with a low fluoride concentration (S1) and stored in the dilute acidic wastewater tank A1. A wastewater from the PSG removal step in the production of crystalline silicon solar cells, which may contain 2˜4% HF, may be collected as an acidic wastewater with a high fluoride concentration (S2) and stored in the concentrated acidic wastewater tank A2. A wastewater from the alkaline texturing step in the production of crystalline silicon solar cells may be collected as an alkaline wastewater (S3) and stored in the alkaline wastewater tank A3. Nowadays, the concentrations of chloride ion in wastewater are regulated in many countries, because high concentrations of chloride are harmful to aquatic plants and animals. To avoid pollution of chloride ion, the traditional and often-used precipitant CaCl 2 may be avoided in a method of the invention. In preferred embodiments, Ca(OH) 2 , CaO, or CaCO 3 may be used. In accordance with embodiments of the invention, in Section B, a new precipitant, an aqueous calcium-containing solution (S4), is prepared by dissolving a known amount of calcium carbonate (alternatively, calcium hydroxide or calcium oxide) stored in the CaCO 3 tank B4, into a certain amount of the collected dilute acidic wastewater (S1) in the stirred dissolution tank B5. The prepared solution (S4) may be filtered and then stored in the Ca 2+ solution tank B6. The prepared solution (S4) may be used as a precipitant in the following processing steps. In addition, to obtain the calcium fluoride products with high purities, the purity of calcium carbonate (calcium hydroxide or calcium oxide) should be high, such as above 80%, preferably over 90%, more preferably over 95%, and most preferably over 98%. CaCO 3 +2H + →Ca 2+ +H 2 O+CO 2 ↓ Ca(OH) 2 +2H + →Ca 2+ +2H 2 O CaO+2H + →Ca 2+ +H 2 O In accordance with embodiments of the invention, recovery zones for concentrated and dilute fluoride-containing wastewaters may be established in the primary treatment unit (Section C) and the secondary treatment unit (Section D), respectively. In Section C (the primary treatment unit), the acidic wastewater with a high fluoride concentration (S2) is passed into the first reactor C7, and the aqueous calcium-containing solution (S4) is also passed into the reactor C7. The concentration of fluoride in the concentrated acidic wastewater (S2) may be measured using spectrophotometric techniques or any other suitable methods. The concentrated acidic wastewater (S2) and the aqueous calcium-containing solution (S4) may be agitated in the reactor C7. The agitation speed may be adjusted to ensure maximum reaction between Ca 2+ ion in the solution S4 and F − ions in the solution S2 and to ensure that all generated calcium fluoride particles are homogeneously distributed throughout the reactor C7. Ca 2+ +2F − →CaF 2 ↑ After completion of the reaction, the reaction slurry in the reactor C7, i.e. the calcium fluoride slurry (S5), is sent from the reactor C7 to a solid-liquid separator C8, wherein the calcium fluoride particles are separated from the wastewater (may be referred to as a mother liquor). The solid-liquid separator C8 may be any suitable equipment known in the art (e.g., centrifuges, filtration equipment, etc.), preferably a centrifuge. After separation, calcium fluoride products (S6) and an effluent from the separator (S7) are obtained. The reaction in the reactor C7 can be carried out at any suitable temperature, for example the temperature of the wastewater. However, the reaction is preferably carried out at a temperature between 30 and 60° C. The residence time for the reaction in the reactor C7 may be from about half an hour to about 3 hours, preferably from about 1 hour to about 2 hours. The concentration of Ca 2+ in the aqueous calcium-containing solution (S4) may be controlled, for example, between about 10 g/l and about 300 g/l, preferably between about 80 g/l and about 120 g/l. The amount of the aqueous solution containing Ca 2+ ion (S4) added to the reactor C7 may be determined by MR1, which is a ratio of the moles of Ca 2+ ion in the solution S4 added into the reactor C7 to the moles of F − ion in the acidic wastewater with high fluoride concentration (S2) added into the reactor C7. MR1 is typically controlled at a value from 0.5 to 1.5, preferably a value of about 0.6. Because two fluoride ions react with one calcium ion, this MR1 ratio is sufficient to ensure most fluoride ion is reacted with calcium ion. The reaction that occurs in reactor C7 normally produces only calcium fluoride particles. This result is due to the high acidity of the acidic wastewater with high fluoride concentration (S2), which is useful in avoiding co-precipitation of silicon and ensuring high purity of the produced calcium fluoride particles. However, the high acidity also results in high solubility of calcium fluoride and therefore a high concentration of fluoride in the effluent from the separator (S7). The effluent from the separator (S7) typically contains a lower fluoride concentration, for example from about 1,000 to about 4,000 mg/l fluoride if the starting wastewater contains more than 20,000 mg/l. The solution S7 may be mixed with the rest of the acidic wastewater with low fluoride concentration (S1) in the mixing tank D9 to get the mixed acidic wastewater (S8) for further treatment. The primary treatment unit results in efficient recovery of a vast majority of fluoride from the acidic wastewater with high fluoride concentration (S2). Typically, if the molar ratio of calcium to fluoride (MR1) is 0.6, 70-90% fluoride in the acidic wastewater with high fluoride concentration (S2) can be recovered. In addition, due to the high acidity (the pH is typically below 3) of the acidic wastewater with high fluoride concentration (S2), the purity of the produced calcium fluoride particles (S6) is above 98%, which is pure enough to be used in many applications, such as for the manufacture of HF. An important thing in Section C (the primary treatment unit) is proper selection of an appropriate reactor type. The selected reactor may be modified to adapt to the specific demands imposed by the characteristics of the wastewater and the specific operational conditions. A stirred tank reactor (STR) may be preferably chosen as the primary treatment unit in accordance with embodiments of this invention for two reasons. First, the acidic wastewater with high fluoride concentration (S2) presents a high acidity (typically with a pH value below 3, such as between 2 and 3), a high fluoride concentration (e.g., ≧2,000 mg/l, preferably ≧5,000 mg/l), and a high ionic strength. Under these conditions, only precipitation process, rather than other methods, allows one to separate and recover fluoride from the wastewater with high efficiencies. Secondly, a stirred tank reactor (STR) is an important precipitation reactor. Industrially, this reactor is easily implemented and easy to control the particle size distribution and the purity of the particles. The process performed in the primary treatment unit (Section C) may be performed in a continuous manner, but preferably in a semi-batch manner because the amount of the acidic wastewater with high fluoride concentration (S2) only accounts for about 1% to 3% of the whole wastewater. Embodiments of the present invention are not limited to the particular methods described here. The secondary treatment unit (Section D) is used to treat the mixed acidic wastewater (S8) which often contains low concentrations of fluoride, such as from about 80 mg/l to 500 mg/l fluoride. The pH value of the mixed acidic wastewater (S8) is relatively low and needs to be raised to at least 6 before the secondary treatment can take place. To raise the pH value, the mixed acidic wastewater (S8) may be mixed with the alkaline wastewater (S3) in the neutralization tank D10 to produce the neutralized wastewater (S9). The pH of the neutralized wastewater (S9) may be controlled between 5 and 8, preferably, between 7 and 8, using an on-line pH meter. The concentration of fluoride in the neutralized wastewater (S9) may be measured using an on-line fluoride ion selective electrode. The aqueous calcium-containing solution (S4), which is prepared in the Section B, may be also used as the precipitator in the secondary treatment unit. In the dilution tank D11, the aqueous calcium-containing solution (S4) may be diluted with the treated wastewater (S11), using a flow calibration tube, such that the concentration of Ca 2+ ions in the diluted calcium reagent solution (S10) is controlled at a suitable level, such as below 10 g/l. The neutralized wastewater (S9) containing fluoride and the diluted calcium reagent solution (S10) are then delivered into the fluidized bed reactor (FBR) D12, in an upward turbulent flow through the supply system and the distributor, in such a way that the solutions are distributed evenly over the cross section of the reactor D12. The reactor D12 contains a fluidized bed of seed materials, which may be the calcium fluoride particles (S6) separated from the solid-liquid separator C8 in Section C. The rates of supply of the solutions S9 and S10 should preferably be controlled to ensure a fluidized state of the bed seed materials. The neutralized wastewater (S9) containing fluoride and the diluted calcium reagent solution (S10) are thoroughly mixed and reacted in the fluidized bed reactor (FBR) D12, wherein crystals of calcium fluoride would form, and the formed calcium fluoride would crystallize on the seed materials. The calcium fluoride seeds grow by a combination of crystal growth and agglomeration between the seed pellets and the newly formed calcium fluoride from the liquid phase. As the seeds grow, they become heavier and tend to move downward within the FBR D12. Over time, the largest particles tend to accumulate in the lower part of the FBR D12, the smallest particles tend to accumulate in the upper part of the FBR D12, and the intermediated-sized particles tend to accumulate in the middle part of the FBR D12. After the FBR D12 has been running for a sufficient duration, the lower part of the FBR D12 will contain pellets of calcium fluoride, which are large enough to be harvested. These large calcium fluoride pellets (S12) may be removed from the reactor D12 periodically, and the fresh seeds are supplemented simultaneously to ensure the maintenance of a well-functioning fluidized bed. After having been removed from FBR D12, the calcium fluoride pellets can be washed and dried. Any suitable drying method may be used. The treated wastewater (S11), which typically contains fluoride at a concentration below 10 mg/l, is then discharged from the FBR D12 into the sewer network of a downstream wastewater treatment plant. If desired, the treated wastewater (S11) may be delivered to a clarifier D13 for reuse or recovery. The reaction in the FBR D12 may be carried out at any suitable temperature, such as the temperature of the wastewater. However, the reaction is preferably carried out at a temperature between 10 and 30° C. As noted above, the calcium fluoride seeds grow through a combination of crystal growth and agglomeration. However, competing reactions such as discrete nucleation and abrasion also take place in the liquid phase. Discrete nucleation and abrasion of the seeds in the liquid phase lead to the formation of fines, which may leave the reactor together with the remaining fluoride in the effluent. The formation of the fines reduces the recovery efficiency of reactor D12. To minimize the formation of fines, the supersaturation in the FBR D12 should be controlled in the metastable region. The control of the supersaturation depends upon the control of the parameters of the wastewater entering the system, such as the temperature, the concentration of the species, and the pH of the wastewater, etc. In accordance with the processes of this invention, the most effective way to maintain the appropriate supersaturation conditions at the inlet of the FBR D12 is to control the concentration of fluoride in the neutralized wastewater (S9) below 200 mg/l, preferably below 150 mg/l. The amount of the diluted calcium reagent solution (S10) added to the FBR D12 is determined by MR2. MR2 is the ratio of the moles of calcium in the solution S10 added into the FBR D12 to the moles of fluoride in the neutralized wastewater (S9) added into the FBR D12. In general, for a fixed concentration of fluoride, the removal ratio of fluoride increases almost linearly with an increase in the MR2, when the MR2 is below 0.5. This is not surprising because two fluoride ions react with one calcium ion to form calcium fluoride. In accordance with embodiments of this invention, MR2 typically has a value of from 0.5 to 2.0, preferably a value of about 0.6-1.5, more preferable from about 0.6-1.0, and most preferably about 0.6. MR2 may be controlled by a molar ratio controller. An MR2 molar ratio controller may comprise any suitable programmable process controller. Those skilled in the field of this invention are familiar with the selection and programming of such controllers. In accordance with embodiments of this invention, for example, an MR2 controller may receive a signal from the fluoride ion selective electrode and a signal from the flow gage as input signals. The fluoride ion selective electrode and the flow gage may be all located at the outlet pipe of the neutralized wastewater feed pump, or any other suitable locations in the flow path. Based on the input signals, the MR2 controller controls a flow control mechanism that determines the addition rate of the diluted calcium reagent solution (S10) into the FBR D12. Any suitable metering mechanism may be used to control the addition rate of the solution (S10) into FBR D12. Such metering mechanisms may include metering pumps, variable valves, or the like. For example, the addition rate of the diluted calcium reagent solution (S10) may be determined as: Q Ca = 2.1 × MR ⁢ ⁢ 2 × Q F × C F C Ca where Q Ca and Q F (L/min) are the influent flows of the diluted calcium reagent solution (S10) and the neutralized wastewater (S9), respectively; C Ca and C F (mg/L) are the concentrations of calcium and fluoride in the diluted calcium reagent solution (S10) and in the neutralized wastewater (S9), respectively; MR2 is the molar ratio as defined above, and the constant 2.1 is the molecular weight ratio of Ca(40)/F(19). The average upward fluid flow velocity, which is measured by dividing the total flow rate of all the streams by the cross sectional area of the reactor D12, should be sufficient to maintain calcium fluoride seeds in suspension. This flow velocity is typically in excess of about 30 m/h. The flow velocity can be controlled by adjusting the feeding pumps of the neutralized wastewater (S9) and the diluted calcium reagent solution (S10) to provide a desired combined flow. The calcium fluoride products produced from Section D (the secondary treatment unit) may contain a low water content, such as about 3-7%. The thus-obtained particles can be used without a drying process for the further industrial applications. The purity of the thus produced calcium fluoride particles (S12) is typically above 97%, which is sufficient for use in the manufacture of HF. The secondary treatment unit may be focused on recovery, less footprint, low reagent cost, safety and easier maintenance. The conditions of supersaturation, pH value, seeds and product quality are the most important factors that need be considered in FBR applications. Changing the operation variables, including the ratio of the amount of calcium reagent and wastewater, fluoride area loading, hydraulic retention times and upflow velocity, could influence the supersaturation in the reactor. It is desirable to use an on-line fluoride ion sensor in the reaction system to control that reaction in a metastable region so as to grow calcium fluoride crystal effectively. With the above processes (Sections A, B, C and D), it is possible to treat fluoride-containing wastewater, such as that produced during crystalline silicon solar cell manufacturing, in an efficient manner. The final treated wastewater streams may have fluoride concentration below 10 mg/L. Particularly, fluoride in the wastewater is mostly recovered as valuable chemical feedstocks, and the amounts of calcium fluoride sludge are reduced to minimal or non-existence. The following example is given to illustrate an embodiment in accordance with the principles of this invention. This example is for illustration only and is not to be viewed as limiting the scope of the present invention. Example A monocrystalline silicon solar cell manufacturer located in Jiangsu province of China was chosen for the demonstration. Section A In this section, a holistic look at the entire manufacture process of the silicon solar cell was taken to identify the various waste streams in the process. Raw wastewater samples were collected on a tank-by-tank basis, and the physico-chemical parameters, such as the pH value, the concentration of fluoride, chemical oxygen demand (COD), total nitrogen (TN), nitrate (N—NO 3 ) and total phosphorus (TP), of the wastewater samples were measured. The analytical data of the major effluents generated at different tanks are presented in Table 1. TABLE 1 The quality parameters of the wastewater generated at different tanks Tanks pH COD N—NH 4 + Total P Total N F − Texturization Tank 1 13.2 488 0.09 17.8 12.9 137 Tank 2 13.8 57000 0.44 17.4 654 104 Tank 3 9.8 <30 0.17 0.05 20.5 2.93 Tank 4 1.2 <30 0.27 <0.01 53.7 0.47 Tank 5 1.7 <30 0.10 0.02 38.5 172 Tank 6 1.2 698 0.08 <0.01 12.7 19800 Tank 7 6.5 <30 0.04 <0.01 23.3 1.19 PSG removal Tank 8 1.3 1430 0.27 <0.01 111 22100 and Edge Tank 9 2.0 <30 0.52 0.30 9.9 156 isolation Tank 10 13.1 138 0.07 63.2 64.1 648 Tank 11 11.1 <30 0.08 0.05 31.9 2.62 In this case, the wastewaters were classified into 3 kinds of streams, according to the concentrations of fluoride and the pH values. The acidic wastewater with a fluoride concentration over 10000 mg/l was collected as a concentrated acidic wastewater (S2), while the acidic wastewater with a fluoride concentration below 200 mg/l was collected as a dilute acidic wastewater (S1). And the wastewater with a pH value over 9 was collected as an alkaline wastewater (S3). The average characteristics of the collected 3 streams are shown in Table 2. TABLE 2 The average characteristics of so-collected 3 kinds of streams Types pH COD N—NH 4 + Total P Total N F − source Alkaline 13.9 25100 0.55 15.9 797 45 Tank 1, 2, 3, 10, 11 wastewater (S3) Concentrated acidic 1.0 462 0.5 0.01 48.2 29200 Tank 6, 8 wastewater (S2) Dilute acidic 1.0 1620 0.35 <0.01 75.3 120 Tank 4, 5, 7, 9 wastewater (S1) Section B In this section, an aqueous calcium-containing solution (S4) is prepared by dissolving a known amount of calcium carbonate into a certain amount of the collected dilute acidic wastewater (S1) in a stirred dissolution tank. The prepared solution is filtered and then stored in the calcium-containing solution tank. The concentration of Ca 2+ in the solution is measured to be 90 g/L. CaCO 3 +2H + →Ca 2+ +H 2 O+CO 2 ↓ Section C The primary treatment were carried out in a laboratory-scale cylindrical PTFE stirred tank crystallizer (1000 ml, d=95 mm) with 4 longitudinal baffles, two inlets for the aqueous calcium-containing solution (S4) and the concentrated acidic wastewater (S2), and a three-blade marine-type impeller covered with polyethylene in order to minimize the secondary nucleation. The crystallizer was immersed in a water bath thermostatically controlled within ±0.5° C. The stirring rate for the experiment was fixed at 300 rpm, which is sufficient to bring about full suspension of the solid-liquid system in the crystallizer. In this example, the experiment was performed in a semi-batch mode. Four hundred (400) mL of the concentrated acidic wastewater (S2) with a fluoride concentration of 29,200 mg/L was poured into the crystallizer and stirred at 300 rpm. When the wastewater was heated to a temperature of 40° C., 136.6 mL of the prior prepared aqueous calcium-containing solution (S4) with a Ca 2+ concentration of 90 g/L was added to the crystallizer at a constant rate of 2 mL/min. After addition of the Ca 2+ aqueous solution, the reaction in the crystallizer was continued for 1 hour. Ca 2+ +2F − →CaF 2 ↓ Then, the calcium fluoride slurry (S5) in the crystallizer was separated by filtration. The separated calcium fluoride particles (S6) was washed and then air-dried at 80° C. for 24 h. The filtrate (S7) was collected for further treatment. In this section, 20.3 g calcium fluoride was recovered from the concentrated acidic wastewater. That is, 84.7% of fluoride in the concentrated acidic wastewater was recovered. The purity of the calcium fluoride was 98.7%. Section D The filtrate (S7) from the section C was mixed with a certain amount of dilute acidic wastewater (S1) to obtain a mixed acidic wastewater (S8). The mixed acidic wastewater (S8) contains about 260 mg/L fluoride. To perform the secondary treatment, the mixed acidic wastewater (S8) is mixed with an alkaline wastewater (S3) to produce a neutralized wastewater (S9). The pH of the neutralized wastewater (S9) was controlled to be around 7. The concentration of fluoride in the neutralized wastewater (S9) was determined to be 110 mg/l. The aqueous calcium-containing solution (S4) prepared in Section B was diluted to obtain a diluted calcium reagent solution (510) containing Ca 2+ 200 mg/L. The secondary treatment was performed in a laboratory-scale fluidized bed reactor (FBR). The reactor was a cylindrical vessel adapted with a clarification compartment according to FIG. 4 . In this example, the various dimensions of the reactor are as follows: D2=2D1=36 mm, and H2=H1/6=63 mm. One skilled in the art would appreciate that this is only an example and that proper dimensions of a reactor may be selected based on the scale of the operation without undue experimentation. The neutralized wastewater (S9) and the diluted calcium reagent solution (S10) with the above mentioned compositions were introduced into the reactor in an upward flow using two peristaltic pumps. The inlet velocities of both solutions were controlled such that the seed materials present in the reactor was carried and maintained in fluidized bed conditions. In this example, the fluid velocity of the neutralized wastewater (S9) was in a range between 0.00015 and 0.0003 m/s, while the fluid velocity of the diluted calcium reagent solution (S10) was between 0.0001 and 0.0003 m/s. The concentration of fluoride in the treated wastewater was analyzed with a fluoride ion electrode (Hach ISEF 121). The temperature in the reaction space was 30° C. In this example, the seed materials were the wet, unwashed calcium fluoride product (S6) produced in Section C, and the amount of the seed materials was 250 g per liter of useful reactor volume. The calcium fluoride thus produced crystallizes out on the surfaces of the seed grains, thereby increasing the particle sizes of seed grains to 100-300 μm. From time to time, these grains were removed from the reactor, and the new seeds were replenished. Due to friction of the calcium fluoride seed grains, calcium fluoride grits (fines) may be formed in the reactor. The grits may be carried along in suspension in the treated wastewater, which is undesirable. In order to prevent or minimize this, a clarification compartment connected to the FBR shown in FIG. 4 was used in this example. The grits thus formed, if any, in the treated water settled in the clarification compartment. Alternatively, a clarifier may be a unit separated from the FBR as shown in FIG. 3 . The operational parameters in the secondary treatment were carefully controlled in this example: The concentration of the fluoride in the neutralized wastewater (S9) was controlled around 110 mg/L; The concentration of Ca 2+ in the diluted calcium reagent solution (S10) was controlled around 200 mg/L; and The molar ratio of Ca 2+ to F − was controlled around 0.8. After passing through the clarifier, the treated wastewater was discharged. The wastewater thus treated has a fluoride content of less than 10 ppm and is free from the said grit. The purity of the calcium fluoride recovered from the treatments was 97%. While the above describe an example in accordance with one embodiment of the invention, one skilled in the art would appreciate that this is only for illustration and not intended to limit the scope of the protection. One skilled in the art after reading this description would appreciate that other modifications and variations are possible without departing from the scope of the invention.

A method for processing fluoride-containing wastewaters from a factory, which includes the following steps: Step 1: collecting the fluoride-containing wastewaters into three pools: an acidic high-fluoride wastewater, an acidic low-fluoride wastewater, and an alkaline wastewater; Step 2: adding a calcium compound to the acidic low-fluoride wastewater to produce a calcium-containing solution; Step 3: reacting a portion of the calcium-containing solution with the acidic high-fluoride wastewater at a calcium-to-fluoride molar ratio of from about 0.5:1 to about 1.5:1 to produce a mixture comprising calcium fluoride particles suspended in a mother liquor; Step 4: separately collecting the calcium fluoride particles and the mother liquor; Step 5: diluting the mother liquor with a diluent to produce a mixed solution; and Step 6: introducing the mixed solution, the calcium-containing solution, and the alkaline wastewater into a fluidized bed reactor, which contains a calcium fluoride crystallization seed material, to form calcium fluoride crystals.

BACKGROUND OF THE INVENTION 1. Field of the Invention This invention relates to a method for the use of a lignan complex isolated from flaxseed for the treatment of atherosclerosis, e.g. reducing or preventing the development of hypercholesterolemic atherosclerosis, for reducing total cholesterol and for raising HDL-C in blood. 2. Description of the Prior Art Hypercholesterolemia is a major risk factor for atherosclerosis (narrowing of the artery due to deposition of fat in the arterial wall) and related occlusive vascular diseases such as heart attack, stroke and other peripheral vascular diseases. Heart disease is the number one killer. Hypercholesterolemic atherosclerosis has been reported to be associated with oxidative stress increase in levels of reactive oxygen species (ROS), production of ROS by polymorphonuclear leukocytes as assessed by chemiluminescence (PMNL-CL), and a decrease in the antioxidant reserve. Pretreatment with antioxidants (vitamin E, probucol, garlic, purpurogallin, secoisolariciresinol diglucoside) reverses the effects of hypercholesterolemia. Flaxseed which is a rich source of α-linolenic acid and richest source of plant lignans has been shown to be effective in reducing hypercholesterolemic atherosclerosis without lowering serum levels of cholesterol. Using flaxseed which has very low α-linolenic acid, has shown that antiatherogenic activity of flaxseed is not due to α-linolenic acid but may be due to lignan component of flaxmeal. Presently the treatment of hypercholesterolemia and hypercholesterolemic atherosclerosis is to reduce hypercholesterolemia by using various lipid lowering agents such as bile acid sequestrant (cholestyramine, colestipol), nicotinic acid, HMG-CoA reductase inhibitor (lovastatin, provastatin, simvastatin, fluvastatin and atrovastatin) and gemfibrozil. Recently probucol which has both antioxidant and lipid lowering activity and vitamin E which has antioxidant activity have been used to prevent atherosclerosis and restenosis. Drugs used for lowering serum lipids and for treatment of atherosclerosis (heart attack and stroke) have many side effects and are expensive. Fibric acid derivatives (gemfibrozil) produces gall stones, myopathy and hepatomegaly. Nicotinic acid produces gastrointestinal symptoms, flushing, hyperglycemia, hepatic dysfunction, elevated uric acid, abnormal glucose tolerance, and skin rash. Bile acid sequestrant (cholestyramine, colestipol) produces gastrointestinal symptoms, and high serum levels of very low density-lipoprotein (VLDL). HMG-CoA reductase inhibitors (statin) produce gastrointestinal symptoms, myopathy and hepatotoxicity. Probucol produces diarrhea and decreases the serum levels of HDL (good cholesterol). Prasad, U.S. Pat. No. 5,846,944, describes the use of secoisolariciresinol diglucoside (SDG), isolated from flaxseed, for reducing hypercholesterolemic atheroscleorsis and reducing serum cholesterol. However, isolating SDG from flaxseed is a relatively expensive procedure. In Westcott et al., U.S. Pat. No. 6,264,853, a new lignan complex is described which has been isolated from flaxseed. This lignan complex typically contains SDG (35%), cinnamic acid glycosides and hydroxymethyl glutaric acid. Only a simple procedure is required to isolate this lignan complex from flaxseed. The purpose of the present invention is to provide a method of using the above lignan complex derived from flaxseed for medical purposes. SUMMARY OF THE INVENTION In accordance with this invention, it has been found that a lignan complex isolated from flaxseed can safely be administered to humans or non-human animals for the treatment of a variety of diseases. The complex and a method for its production are described in Westcott et al., U.S. Pat. No. 6,264,853, issued Jul. 24, 2001, and incorporated herein by reference. This complex is used in substantially pure form, e.g. a purity of at least 95%, and contains secoisolariciresinol diglucoside (SDG), cinnamic acid glucosides and hydroxymethyl glutaric acid. It typically has a nominal molecular weight in the range of about 30,000 to 100,000. The complex can be administered orally or intraperitoneally and has been found to be highly effective when administrated in a daily oral dosage of 20 to 60 mg per kg of body weight. The oral doses may conveniently be in the form of tablets or capsules and the complex may be used together with a variety of pharmaceutically acceptable diluents or carriers. When administered to humans or non-human animals, the complex has been found to be highly effective for treating hypercholesterolemic atherosclerosis, as well as for reducing total cholesterol and raising HDL-C in blood. Thus, it is useful for the prevention and treatment of coronary artery disease, stroke and other peripheral vascular diseases. BRIEF DESCRIPTION OF THE DRAWINGS In the drawings which illustrate the present invention: FIG. 1 is a graph showing sequential changes in serum triglyceride concentration for four different experimental groups; FIG. 2 is a graph showing sequential changes in serum total cholesterol concentration of four different experimental groups; FIG. 3 is a graph showing sequential changes in serum LDL-C concentration for four different experimental groups; FIG. 4 is a graph showing sequential changes in serum HDL-C concentration for four different experimental groups; FIG. 5 is a bar graph showing sequential changes in serum malondialdehyde (MDA) for four different experimental groups; FIG. 6 is a bar graph showing aortic tissue malondialdehyde (MDA) concentration for four different experimental groups; FIG. 7 is a bar graph showing aortic tissue chemiluminescence (Aortic-CL) for three different experimental groups; FIG. 8 shows photographs of endothelial surfaces of aortas for four different experimental groups; and FIG. 9 is a bar graph showing the extent of atherosclerotic plaques in the initial surface of aorta for four different experimental groups. In the graphs, the results are expressed as mean±S.E. The symbols used in FIGS. 1 to 5 have the following meanings. *P<0.05, Comparison of values at different times with respect to time “0” in the respective group. a P<0.05, Control vs other groups. b P<0.05, Lignan complex vs 0.5% cholesterol or 0.5% cholesterol+lignan complex. c P<0.5, 0.5% Cholesterol vs 0.5% cholesterol+lignan complex. + P<0.05, Month 1 vs month 2 in the respective groups. In FIGS. 6 and 7 the symbols have the following meanings: *P<0.05, control Vs other groups; †P<0.05, lignan complex Vs 0.5% cholesterol or 0.5% cholesterol+lignan complex. a P<0.05, 0.5% cholesterol Vs 0.5% cholesterol+lignan complex. For FIG. 7, the significance symbols are as follows: *P<0.05, control Vs other groups. †P<0.05, 0.5% cholesterol Vs 0.5% cholesterol+lignan complex. DESCRIPTION OF THE PREFERRED EMBODIMENTS The complex used according to this invention typically contains about 34 to 37% by weight of SDG, about 15 to 21% by weight cinnamic acid glucosides and about 9.6 to 11.0% by weight hydroxmethyl glutaric acid. The cinnamic acid glycosides include coumaric acid glucoside and ferulic acid glucoside. They are typically present in the complex in amounts of about 9.5 to 16.0% by weight coumaric acid glucoside and 4.5 to 5.0% ferulic acid glucoside. The complex composition typically contains about 59 to 70% by weight of the above active ingredients. The balance comprises protein, ash and water of crystallization. The invention is further illustrated by the following non-limiting examples. EXAMPLE 1 EXPERIMENTAL PROTOCOL Experiments were conducted on New Zealand White rabbits. Rabbits were assigned to 4 groups as shown in Table 1. Those in Group I were fed rabbit laboratory chow diet. The other groups received lignan complex or cholesterol or cholesterol+lignan complex. The lignan complex was obtained from Agriculture and Agri-Food Canada and was extracted from flaxseed by the method described in Westcott et al., U.S. Pat. No. 6,264,853incorporated herein by reference. The diet was especially prepared by Purina and did not contain any antioxidant. Lignan complex was given orally daily in the dose of 40 mg/kg body weight. The rabbits were cared for according to approved standards for laboratory animal care. The rabbits were on their respective diet treatment for 2 months. Blood samples were collected (from ear marginal artery) for measurement of serum-triglycerides (TG), total cholesterol. (C), low-density lipoprotein cholesterol (LDL-C), high-density lipoprotein cholesterol (HDL-C), enzymes, albumin, creatinine, and malondialdehyde (MDA) before (0 time) and after 1 and 2 months on the respective experimental diets. The rabbits were anesthetized at the end of 2 months and aortas were removed for assessment of atherosclerotic plaques, and measurement of aortic tissue MDA and antioxidant reserve (Aortic-chemiluminescence). The measurement of lipids, atherosclerotic plaques, oxidative stress were done according to known methods. Serum enzymes, albumin and creatinine for assessment of liver and kidney function were measured by already established techniques. Assessment of hemopoietic system were made by established techniques available in the hospital. SERUM LIPIDS. Changes in serum TG, C, LDL-C, and HDL-C in the 4 groups are shown in FIGS. 1-4. Lignan complex did not affect serum TG, TC, LDL-C but increased HDL-C significantly in the groups on control diet. A 0.5% cholesterol diet increased serum TG, C, LDL-C and HDL-C. Lignan complex in 0.5% cholesterol-fed rabbit produced less increase in C and LDL-C, and greater increase in HDL-C as compared to only 0.5% cholesterol-fed rabbits. Serum TG levels were similar in Group III and IV. These results indicate that the lignan complex lowers serum cholesterol (significantly) and LDL-C (not significant), and raises HDL-C (significantly) in hypercholesterolemic rabbit. Lignan complex also raises HDL-C in normocholesterolemic rabbits. OXIDATIVE STRESS. Results for oxidative stress parameters (serum MDA, aortic tissue-MDA, aortic tissue antioxidant reserve) are shown in FIGS. 5-7. Serum MDA levels remained unaltered in control and lignan complex groups. It increased in both 0.5% cholesterol and 0.5% cholesterol+lignan groups. However, the increase was less in groups with 0.5% cholesterol+lignan complex. Aortic MDA increased and lignan complex decreased in 0.5% cholesterol-fed rabbits. Aortic tissue chemiluminscence (Aortic-CL) is a measure of antioxidant reserve. An increase in Aortic-CL suggests a decrease in the antioxidant reserve and vice-versa. Aortic-CL decreased in cholesterol-fed group of rabbits. Lignan complex in cholesterol-fed rabbits tended to increase the aortic-CL compared to 0.5% cholesterol without lignan complex. These results indicate that high cholesterol increases oxidative and the lignan complex reduces oxidative stress. ATHEROSCLEROSIS. Representative photographs of endothelial surfaces of aortas from each group are depicted in FIG. 8, and the results are summarized in FIG. 9 . In FIG. 8, Group I is Control, Group II is lignan complex, Group III is 0.5% cholesterol and Group IV is 0.5% cholesterol+lignan complex. In FIG. 9 : *P<0.05Group I or Group II vs Group III and Group IV. †P<0.05, Group III vs Group IV. Atherosclerotic plaques were absent in Group I and II. However, a significant area of aortic surface from Group III (50.84±6.23%) and Group IV (33.40±4.80%) was covered with atherosclerotic plaques. This indicates that the lignan complex reduced the hypercholesterolemic atherosclerosis by 34.3%. HEMOPOIETIC SYSTEM Red Blood Cells (RBCs). The changes in various parameters related to RBC are shown in Tables 2-8. In general lignan complex in the control diet group (Group II) did not affect the RBC count, hemoglobin (Hb), hematocrit (Hct), mean corpuscular volume (MCV), mean corpuscular hemoglobin (MCH), mean corpuscular hemoglobin concentration (MCHC) and red blood cell distribution width (RDW). Cholesterol diet (Group III) alone produced significant decreases in RBC, Hb, Hct and MCH; increases in RDW; and no change in MCV and MCHC. Lignan complex in 0.5% cholesterol-fed rabbits (Group IV) reduced RBC, Hb, and Hct; increased MCV, MCH and RDW. The values for RBC, Hb, Hct, MCV, MCH, MCHC and RDW in Group IV were not significantly different from those in Group III. This shows that, in general, the lignan complex has no adverse effects on the hemopoietic system. White Blood Cells. The changes in the white blood cells (WBCs) and the differential counts granulocytes, lymphocytes, and monocytes are shown in Tables 9-12. Lignan complex in the control diet group (Group II) produced decreases in WBCs and monocytes, and no changes in granulocytes and lymphocytes. These changes in the various parameters in Group II were not significantly different from those in control group (Group I). These parameters of WBCs were unaffected in Group III and IV except in Group III where monocyte counts decreased. These results indicate that lignan complex has no adverse effects on the WBCs, granulocytes, lymphocytes and monocyte counts. PLATELET. The changes in platelet counts and mean platelet volume (MPV) of the four groups are summarized in Tables 13-14. Platelet counts slightly decreased in Group I but MPV remained unchanged. These parameters remained unaltered in Group II. Basically, all the parameters in all the groups remained unaltered. These results indicate that lignan complex has no adverse effects on platelet counts and mean platelet volume. EXAMPLE 2 Studies were conducted to determine if the lignan complex given for 2 months produces adverse effects on liver and kidney function. (a) Assessment of liver function was made by measuring serum enzymes [alkaline phosphatase (ALP), alanine amino-transferase (ALT), aspartate aminotransferase (AST) and gamma-glutamyltransferase (GGT)] and serum albumin. These serum enzymes are elevated and serum albumin is decreased in liver disease. The results are summarized in Table 15-19. Serum levels of ALT, AST and GGT were similar in Groups I and II at month two of the protocol, however levels of serum ALP were lower in Group II compared to Group I. The changes in the serum levels of ALP, ALT and GGT remained unchanged as compared to “0” month in the Groups III and IV. However serum levels of AST increased to a similar extent in both groups III and IV. Serum albumin levels increased at month one as compared to “0” month in all the groups, however the increases at month two were not significantly different as compared to “0” month. The values of serum albumin at month two, although higher in Groups I and II as compared to Group III and IV, they were not significantly different from each other. These results indicate that hypercholesterolemia has adverse effects on liver function and that the lignan complex does not have adverse effects on liver function. (b) Assessment of kidney function was made by measuring serum enzymes (ALT and AST) and creatinine. ALT, AST and creatinine levels are elevated in dysfunctional kidney. The results are summarized in Tables 16, 17 and 20. There were no significant differences in the values of serum ALT, AST and creatinine among the 4 groups. These results indicate that the lignan complex or hypercholesterolemia did not have adverse effects on kidney function. EXAMPLE 3 The lignan complex was also fed orally to normal ratsfor 2 months at a daily dosage of 40 mg/kg of body weight and the rats were studied to see if the complex had any affect on the liver and kidney function and hemopoietic cells. It was found that the lignan complex did not affect any of the above, indicating that it is not toxic to liver, kidney and blood cells. TABLE 1 Experimental Diet Groups Group Diet/Treatment I (n = 10) Control (Rabbit chow diet) II (n = 6) Lignan complex control (Rabbit chow diet supplemented with lignan complex, 40 mg/kg body weight, orally, daily) III (n = 12) Cholesterol diet (0.5% cholesterol in rabbit chow diet) IV (=16) Cholesterol diet + lignan complex (0.5% cholesterol diet supplemented with lignan complex, 40 mg/kg; body weight, orally, daily) TABLE 2 Red Blood Cells (RBC) Counts (10 12 /L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 5.08 ± 0.12 6.22 ± 0.18* 6.04 ± 0.16* II. Control diet + 5.76 ± 0.16 a 6.03 ± 0.13 5.90 ± 0.20 lignan complex III. 0.5% cholesterol 5.67 ± 0.07 a 5.61 ± 0.09 a,b 4.60 ± 0.11* ,†,a,b diet IV. 0.5% cholesterol 5.83 ± 0.09 a 5.43 ± 0.06* ,a,b 4.70 ± 0.14* ,†,a,b diet + lignan complex *P < 0.05, 0 month vs 1 and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III or Group IV. TABLE 3 Hemoglobin Levels in the Blood (g/L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 111.8 ± 1.6 133.8 ± 4.2* 128.9 ± 2.5* II. Control diet + 126.8 ± 2.2 † 134.6 ± 2.4* 129.0 ± 3.1 lignan complex III. 0.5% cholesterol 125.1 ± 1.6 † 122.8 ± 1.7 †,a 107.7 ± 2.2* ,†,a diet IV. 0.5% cholesterol 127.6 ± 1.6 † 122.9 ± 1.8 †,a 110.0 ± 2.6* ,†,a diet + lignan complex *P < 0.05, 0 month vs 1 and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III or Group IV. TABLE 4 Hematocrit (L/L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 0.327 ± 0.00 0.385 ± 0.01* 0.376 ± 0.01* II. Control diet + 0.357 ± 0.00 a 0.382 ± 0.01* 0.373 + 0.01 lignan complex III. 0.5% choles- 0.364 ± 0.01 a 0.348 ± 0.01* ,a,b 0.300 ± 0.01* ,†,a,b terol diet IV. 0.5% choles- 0.372 ± 0.00 a 0.347 ± 0.01* ,a,b 0.310 ± 0.01* ,†,a,b terol diet + lignan complex *P < 0.05, 0 month vs 1 and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III or Group IV. TABLE 5 Mean corpuscular volume (fL) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 64.4 ± 0.9 61.8 ± 0.9 62.3 ± 0.9 II. Control diet + 62.4 ± 1.1 63.4 ± 1.0 63.3 ± 0.8 lignan complex III. 0.5% cholesterol 64.2 ± 0.4 62.2 ± 0.4* 65.3 ± 0.5 †,a diet IV. 0.5% cholesterol 63.8 ± 0.6 64.0 ± 0.6 c 66.1 ± 0.6* ,†,a,b diet + lignan complex *P < 0.05, 0 month vs 1 and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III or Group IV. c P < 0.05, Group III vs Group IV. TABLE 6 Mean Corpuscular Hemoglobin (pg) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 22.0 ± 0.34 21.5 ± 0.26 21.4 ± 0.32 II. Control diet + 22.1 ± 0.38 22.4 ± 0.33 21.9 ± 0.34 lignan complex III. 0.5% cholesterol 22.1 ± 0.15 22.0 ± 0.24 23.4 ± 0.20* ,†,a,b diet IV. 0.5% cholesterol 21.9 ± 0.26 22.7 ± 0.19* ,a,c 23.4 ± 0.23* ,†,a,b diet + lignan complex *P < 0.05, 0 month vs 1 and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III or Group IV. c P < 0.05, Group III vs Group IV. TABLE 7 Mean Corpuscular Hemoglobin Concentration (g/L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 341.8 ± 3.8 347.6 ± 1.8 343.3 ± 1.5 II. Control diet + 354.2 ± 1.7 352.5 ± 1.7 346.2 ± 2.4* lignan complex III. 0.5% cholesterol 343.5 ± 1.5 351.7 ± 2.5 357.7 ± 1.9 diet IV. 0.5% cholesterol 342.9 ± 2 354.6 ± 3.0 354.9 ± 1.7 diet + lignan complex *P < 0.05, 0 month vs 1 and 2 months in the respective groups. TABLE 8 Red Blood Cell Distribution Width (RDW) as % in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 11.68 ± 0.22 12.44 ± 0.37 12.84 ± 0.35 II. Control diet + 13.52 ± 0.28 a  13.1 ± 0.33 12.95 ± 0.32 lignan complex III. 0.5% cholesterol 11.56 ± 0.20 b  12.3 ± 0.14* ,b 13.42 ± 0.3* ,† diet IV. 0.5% cholesterol  12.2 ± 0.27 b  13.1 ± 0.29* ,c 13.17 ± 0.21* diet + lignan complex *P < 0.05, 0 month vs 1 and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III or Group IV. c P < 0.05, Group III vs Group IV. TABLE 9 White Blood Cell Counts (10 9 /L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 4.96 ± 0.62 5.34 ± 0.77 5.3 ± 0.38 II. Control diet +  7.7 ± 0.7 a 7.33 ± 0.25 a 4.85 ± 0.78* ,† lignan complex III. 0.5% cholesterol 6.34 ± 0.33  8.5 ± 0.48* ,a 6.88 ± 0.73 diet IV. 0.5% cholesterol 6.05 ± 0.28 b 9.14 ± 0.44* ,a,b 6.59 ± 0.93 † diet + lignan complex *P < 0.05, 0 month vs 1 and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III or Group IV. TABLE 10 Granulocytes Content of Blood (10 9 /L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 0.96 ± 0.12 0.64 ± 0.08 1.02 ± 0.15 II. Control diet + 1.14 ± 0.18 1.21 ± 0.08 a 0.65 ± 0.09 † lignan complex III. 0.5% cholesterol 1.10 ± 0.08 1.75 ± 0.27* ,a 1.36 ± 0.23 diet IV. 0.5% cholesterol 0.87 ± 0.078 1.21 ± 0.20 1.84 ± 0.38* diet + lignan complex *P < 0.05, “0” time vs 1 month and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. TABLE 11 Lymphocyte Counts in Blood (10 9 /L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 3.42 ± 0.44 4.22 ± 0.69  4.0 ± 0.29 II. Control diet + 5.27 ± 0.33 a 5.53 ± 0.27 a 3.87 ± 0.65 † lignan complex III. 0.5% cholesterol 4.19 ± 0.19 b 6.41 ± 0.51* 4.74 ± 0.46 †t diet IV. 0.5% cholesterol 4.65 ± 0.27 a 5.02 ± 0.74 5.09 ± 0.56 diet + lignan complex *P < 0.05, 0 month vs 1 and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III or Group IV. TABLE 12 Monocyte Counts in the Blood (10 9 /L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 0.56 ± 0.07 0.46 ± 0.05 0.28 ± 0.07 * II. Control diet + 0.62 ± 0.07 0.58 ± 0.047 0.33 ± 0.08 *,† lignan complex III. 0.5% cholesterol 0.75 ± 0.07 0.75 ± 0.08 a 0.37 ± 0.05 *,† diet IV. 0.5% cholesterol 0.45 ± 0.05 c 0.56 ± 0.09 0.45 ± 0.02 a diet + lignan complex *P < 0.05, “0” month vs 1 month and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. c P < 0.05, Group III or Group IV. TABLE 13 Platelet Counts in the Blood (10 9 /L) in the Blood of Various Experimental Groups Time (months) Group 0 1 2 I. Control diet 393 ± 46 324 ± 52 286 ± 25* II. Control diet + 329 ± 20 280 ± 8* 267 ± 23 lignan complex III. 0.5% cholesterol 422 ± 24 b 341 ± 26* 401 ± 31 a,b diet IV. 0.5% cholesterol 403 ± 20 b 309 ± 23* 364 ± 37 diet + lignan complex *P < 0.05, “0” month vs other months in the respective groups. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III or Group IV. TABLE 14 Mean Platelet Volume in Fentoliter (f1) for Various Experimental Groups Time (months) Group 0 1 2 I. Control diet 5.38 ± 0.27 5.64 ± 0.21 5.69 ± 0.21 II. Control diet + lignan complex 6.03 ± 0.16 5.81 ± 0.11 5.75 ± 0.11 III. 0.5% cholesterol diet 5.37 ± 0.07 a 5.27 ± 0.09 a 5.96 ± 0.11 *,† IV. 0.5% cholesterol diet + lignan complex 5.51 ± 0.7 a 5.47 ± 0.07 a 5.85 ± 0.11 *,† * P < 0.05, “0” month vs other months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective groups. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III or Group IV. TABLE 15 Serum Alkaline Phosphatase (ALP) Levels (U/L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 121.3 ± 20.2 152.6 ± 10.9 118.1 ± 7.7 † II. Control diet + 71.0 ± 9.65 a lignan complex III. 0.5% cholesterol 169.1 ± 16.6 191.5 ± 15.6 160.3 ± 11.6 a,b diet IV. 0.5% cholesterol 142.7 ± 1.4 181.2 ± 7.9* 132.7 ± 19.5 diet + lignan complex *P < 0.05, “0” month vs 1 month and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. b P < 0.05, Group II vs Group III and Group IV. TABLE 16 Serum Alanine Aminotransferase (ALT) Levels (U/L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 27.25 ± 3.6 44.2 ± 5.85  41.2 ± 3.5* II. Control diet + Not Not 47.33 ± 9.2 lignan complex measured measured III. 0.5% cholesterol 41.22 ± 3.1 a 59.3 ± 11.2 69.08 ± 13.3 diet IV. 0.5% cholesterol  41.6 ± 2.6 a 45.4 ± 9.4  39.1 ± 5.6 diet + lignan complex *P < 0.05, “0” month vs other months in the respective groups. a P < 0.05, Group I vs other groups. TABLE 17 Serum Aspartate Aminotransferase (AST) Levels (U/L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 25.0 ± 5.1 25.4 ± 0.5 41.1 ± 3.6 *,† II. Control diet + Not Not 34.7 ± 4.4 lignan complex measured measured III. 0.5% cholesterol 35.0 ± 3.3 44.1 ± 6.5 53.1 ± 2.4 *,a diet IV. 0.5% cholesterol 28.8 ± 4.4 29.6 ± 2.4 49.4 ± 4.5 *,† diet + lignan complex *P < 0.05, “0” month vs 1 month and 2 months in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. TABLE 18 Serum Levels (U/L) of Gamma-glutamyltransferase (GGT) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 9.0 ± 0.8 8.8 ± 1.0 8.0 ± 1.9 II. Control diet + Not Not 8.0 ± 1.3 lignan complex measured measured III. 0.5% cholesterol 9.6 ± 0.4 8.6 ± 0.8 6.4 ± 1.7 diet IV. 0.5% cholesterol 9.0 ± 1.1 6.5 ± 0.6 6.4 ± 1.2 diet + lignan complex TABLE 19 Serum Albumin Levels (gm/L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet 15.8 ± 0.2 17.8 ± 0.37*  30.3 ± 5.2 II. Control diet + 35.50 ± 4.52 lignan complex III. 0.5% cholesterol 16.9 ± 0.31 18.3 ± 0.42* 20.41 ± 2.58 diet IV. 0.5% cholesterol 17.2 ± 0.2 19.0 ± 0.54* 26.77 ± 4.1 diet + lignan complex *P < 0.05, comparison of the values at various times with respect to “0” time in the respective groups. TABLE 20 Serum Creatinine Levels (μmoles/L) in the Experimental Groups Time (months) Group 0 1 2 I. Control diet  48.4 ± 2.46 78.8 ± 2.35 *  103.3 ± 5.4 *,† II. Control diet + Not Not 104.25 ± 6.2 * lignan complex measured measured III. 0.5% cholesterol 62.12 ± 1.68 a 80.7 ± 3.9 *  106.4 ± 4.6 *,† diet IV. 0.5% cholesterol  60.0 ± 6.0 73.2 ± 3.7  97.56 ± 5.33 *,† diet + lignan complex * P < 0.05, comparison of values at various times with respect to “0” time in the respective groups. † P < 0.05, 1 month vs 2 months in the respective group. a P < 0.05, Group I vs other groups. Since lignan complex lowers serum cholesterol, elevates serum HDL-C and reduces hypercholesterolemic atherosclerosis it will be of use in the prevention and treatment of the following diseases: i) Hypercholesterolemic atherosclerosis. ii) Coronary artery disease (heart attack). iii) Stroke. iv) Restenosis following percutaneous transluminal coronary angioplasty. v) Restenosis after stent implant. vi) Stroke, heart attack, renal failure and retinopathy in diabetes mellitus. vii) Hypercholesterolemia. viii) Peripheral vascular diseases, such as intermittent clandication. The use of lignan complex derived from flaxseed according to this invention has the following advantages: i) Lignan complex contains materials that have antioxidant and anti-PAF activity and hence is an anti-inflammatory agent. ii) It lowers serum cholesterol, raises HDL-C and reduces hypercholesterolemic atherosclerosis. iii) This compound is a natural food product and has no toxicity on hemopoietic system, liver and kidney, and it is a safe drug. iv) It is inexpensive and safe as compared to other drugs used for lowering lipids and reducing atherosclerosis. v) This compound is cheaper than SDG because processing of SDG is expensive as compared to lignan complex. vi) The dose of lignan complex is very small as compared to flaxseed.

A method is described for treating hypercholesterolemic atherosclerosis or for reducing total cholesterol while raising high-density lipoportoein cholesterol. It involves administering to a patient a substantially pure complex derived from flaxseed and containing secoisolariciresinol diglucoside (SDG), cinnamic acid glucosides and hydroxymethyl glutaric acid.

[0001] This application claims priority of U.S. patent application Ser. No. 13/852,614 filed on Mar. 28, 2013 which application claims priority of U.S. patent application Ser. No. 13/645,184 filed on Oct. 4, 2012 which claims priority of U.S. Provisional Patent Application Ser. No. 61/544,468, entitled “APPARATUS AND METHOD FOR THE CONDENSED PHASE PRODUCTION OF TRISILYLAMINE” filed on Oct. 7, 2011, which is incorporated by reference herein. FIELD OF THE INVENTION [0002] This invention relates to a batch method for synthesizing silylamines, particularly trisilylamine in a solvent. The invention relates to a process that promotes reaction conditions suitable for a high efficiency synthesis of silylamines. The primary silylamine of interest is trisilylamine. Production of disilylamine in commercial quantities is also within the scope of the present invention. BACKGROUND OF THE INVENTION [0003] Trisilylamine (“TSA”) is a useful molecule for use in semiconductor manufacturing. It is stable once produced, but is susceptible to decomposition from excessive reaction conditions and synthesis by-products. Dussarrat, et al. U.S. Pat. No. 7,192,626 demonstrated that a stable Silicon nitride film is formed on a substrate by feeding trisilylamine and ammonia into a CVD reaction chamber that contains a substrate. [0004] Wells and Schaeffer ( J. Am. Chem. Soc., 88:1, 37 (1996)) discuss a batch method of preparing trisilylamine by the reaction silyl chloride with ammonia. They report the yield of trisilylamine varied depending on the method of mixing and the purity of the reactants. Wells and Schaeffer allowed the reactants to mix in the gas phase by introducing the ammonia from below into a 1 liter bulb containing silylchloride. After introducing the gaseous ammonia very slowly, the reaction bulb and contents were allowed to remain at room temperature for 15 min. Copious amounts of white solid were precipitated on the walls of the bulb as soon as mixing occurred. The product was removed and the trisilylamine recovered. The process yield was about 77% of the theoretical amount of trisilylamine. [0005] In the batch reactor process, all of monohalosilane is charged into the reactor vessel. Batch size is limited by this initial charge and the size of the vessel. Ammonia gas is then slowly added into the flask. Reaction conditions will vary in the vessel depending on the initial concentrations of monohalosilane and ammonia and the efficiency of turbulent mixing in the vessel. The mixing is affected by vessel size as well as the efficiency of the mechanical mixing device if one is employed. In addition, during the batch process the silylamines produced are in contact with ammonium halide which is also a product of the reaction. Ammonium halides such as ammonium chloride are catalysts and will disproportionate TSA into silane and other degradation products thereby lowering the yield of TSA. The reaction of silyl halide and ammonia produces heat thereby exacerbating the degradation conditions in a closed reactor vessel. [0006] US 2010/0310443 is directed to a tubular flow gas phase reactor and a process for the synthesis of silylamines which have been found to produce high volumes, at high yield efficiencies of silylamines. The reactor has a combination of characteristics found in plug flow and laminar flow devices. This combination of properties results in a high volume high efficiency synthesis of silylamines. The primary silylamine of interest is trisilylamine. Production of disilylamines in commercial quantities is also within the scope of the present invention. This process produces high volumes of ammonium halide requiring the reaction tube to be opened and cleaned after each production batch is produced. This is a labor intensive process leading to significant down time. SUMMARY OF THE INVENTION [0007] The present invention is directed to a condensed phase batch process for synthesis of TSA comprising: (a) adding a solvent to a reactor vessel; (b) cooling the solvent; (c) condensing monohalosilane into the solvent to form a solution; (d) adding anhydrous ammonia into the solution to form a reaction mixture; (e) separating the silylamines, excess monohalosilane and TSA from the reaction mixture; and (f) purifying the silylamines to obtain TSA; [0008] Condensed phase reactions of excess monohalosilanes, such as monochlorosilane (“MCS”) with ammonia, are beneficial since the formation of TSA occurs rapidly concomitantly producing ammonium halide salt which in the case of MCS is ammonium chloride (“NH 4 Cl”) salt. Such salts are localized in the reaction zone as a slurry with the solvent. This approach preferably utilizes a high boiling point solvent to act as a heat-transfer medium in which the ammonium chloride salt is dispersed and downstream product removal is devoid of salt formation. The general benefit of this approach is the formation of TSA in the condensed phase followed by vacuum stripping of the product from the reaction slurry and discharge of the waste salt/solvent slurry from the reactor vessel after which the reactor can be re-charged with solvent and excess, liquefied monohalosilane for another batch synthesis. In this condensed phase process, the reactor does not have to be cleaned before the next batch run as the ammonium chloride salt byproduct of the reaction is removed as a slurry in the solvent. [0009] This condensed phase reaction scheme provides the following benefits: [0010] A. Low temperature, condensed-phase reactions of ammonia with monohalosilane in a solvent in which the formation of TSA is enhanced over a relatively short period of time. [0011] B. Suitable solvents such as anisole (methoxybenzene) provide vapor pressure depression/boiling point elevation of the MCS reagent, which promotes the formation of liquefied MCS and favorable condensed-phase disilylamine (“DSA”) intermediate reaction kinetics. [0012] C. The solvent acts as a uniform heat transfer medium in which byproduct waste salt is dispersed and localized predominantly in the reaction mixture. [0013] D. The suppression of partially substituted silylamines (such as DSA) that could react further down stream during product collection. [0014] The complete reaction is: [0000] 4NH 3 +3SiH 3 X→3NH 4 X+(SiH 3 ) 3 N [0015] It is believed that the silylamines of the present invention are produced in accord with the following reaction sequence: [0000] 2NH 3 +SiH 3 X→NH 4 X+SiH 3 NH 2 [0000] 2SiH 3 NH 2 →NH 3 +(SiH 3 ) 2 NH [0000] 3(SiH 3 ) 2 NH→NH 3 +2(SiH 3 ) 3 N [0016] Where X=Cl, F, Br, I BRIEF DESCRIPTION OF THE DRAWINGS [0017] Specific embodiments of the invention are described below with reference to the following. [0018] FIG. 1 is a simplified schematic diagram of a condensed phase reactor of the invention utilizing a Schlenk tube. [0019] FIG. 2 is a simplified schematic diagram of a condensed phase reactor of the invention utilizing a Parr reactor vessel. [0020] FIG. 3 is the graphic representation of temperature and pressure versus time for Experiment 8. [0021] FIG. 4 is the graphic representation of temperature and pressure versus time for Experiment 9. [0022] FIG. 5 is the graphic representation of temperature and pressure versus time for Experiment 11. [0023] FIG. 6 is the graphic representation of temperature and pressure versus time for Experiment 12. [0024] FIG. 7 is the graphic representation of temperature and pressure versus time for Experiment 13. [0025] FIG. 8 is the graphic representation of solvent polarity (ET N ) and Lewis Basicity (Donor Number “DN”) for specific solvents. DETAILED DESCRIPTION OF THE INVENTION [0026] The general method of this invention includes the following: [0027] Filling the reactor with an appropriate solvent (anisole, high boiling point ethers, aliphatic and aromatic hydrocarbons, etc.); [0028] Adjusting the solvent or solvent formulation (one or more solvents) to an initial temperature between about 100° C. and about −78° C.; [0029] Adding an excess of monohalosilane (relative to the final amount of ammonia added to the system) in a range of about 5 to about 300 mole %; [0030] Adding anhydrous ammonia into the solution through one or more dip/sparge tubes that are below the liquid level at a rate that silylamine formation and does not induce the formation of polysilazane and silane and reactions of the solvent with the reactants to form undesired products. Forming silylamines in the reaction mixture; A non-limiting list of factors affecting the rate of addition of ammonia into the solution include, volume of the solvent, concentration of the monohalosilane, temperature of the solvent reaction mixture, mixing efficiency, and the rate of heat transfer out of the reaction vessel. A preferred rate of addition of ammonia for moderate sized batch reactions would be from about 100 mg to 5 g/minute, for larger and production batches the rate of addition would be a function of batch size and therefore would be proportionately greater; a preferred temperature of the reaction solvent throughout the addition of ammonia is from about 70° C. to just above the freezing point of the solvent and reactant solution. [0031] Separating the reaction products, silylamines, from the reaction solution. Reaction product silylamines are trisilylamine and disilylamine. Preferred methods of separation include vacuum stripping or distilling the product mixture, which may be preceded by filtration, at reduced pressure once all of the ammonia has been added and collecting the distillates which contain the product(s) in a low temperature cryotrap. The temperature of the reaction mixture may be raised during vacuum stripping. In pilot scale batches the reaction mixture temperature has be raised to about 100° C. during vacuum stripping. [0032] Purifying the product to obtain pure aminosilane. The preferred aminosilane is trisilylaming (“TSA”); preferred purification processes are fractionation or distillation. [0033] Mixing the waste ammonium chloride salts in solvent to suspend or slurry the solids in the solvent and either drain or pressure transfer the waste stream out of the reactor. [0034] The reactor can then be re-charged for another batch synthesis). [0035] A step wise description of the process follows: [0036] A process for preparing silylamine comprising: (a) adding a solvent to a reactor vessel; (b) adding monohalosilane into the solvent to form a solution; (c) adding anhydrous ammonia into the solution to form a reaction mixture; (d) forming silylamines in the reaction mixture; (e) separating silylamines from the reaction mixture; and (f) purifying silylamines; [0043] wherein the solvent has a DN between about 6 to about 28 and an E T N from about 0.1 to about 0.4. [0044] In a preferred embodiment, anisole is the solvent and an excess of about 20 to about 50 mole % monohalosilane to ammonia is used and an initial reaction temperature of about 10° C. to about 60° C. In a preferred embodiment, the monohalosilane is MCS. [0045] The preferred ammonia addition process is to react the ammonia and MCS in the solution and limit any gas-phase reactions in the headspace above the solution thereby avoiding ammonium chloride build up on the exposed surface of the reactor vessel and down stream of the reaction vessel such as in the cryogenic traps. Ammonium chloride salt found in the downstream storage vessels is referred to as “down stream salt.” [0046] The solvent aspect ratio is defined as the relationship of the height of solvent (liquid level) divided by the internal diameter of the reactor and is important relative to the path the ammonia or intermediate disilylamine (DSA) product has to travel to break through the surface of the liquid at the solvent-headspace interface. [0047] The lower limit value for the aspect ratio is not critical but is an experienced based guide for setting an anhydrous ammonia gas flow/feed rate in a particular reactor in light of the following parameters; solvent, MCS concentration, temperature and pressure. [0048] Preferred operation of the reactor is achieved when the feed rate of the ammonia gas is adjusted such that all of the ammonia is reacted with MCS in solution and none of ammonia gas escapes the solution to enter the headspace above the solvent surface. Better gas dispersion methods, better mixing and a higher solvent aspect ratio are process methods that will support a higher ammonia gas flow rate thereby speeding processing time. [0049] The preferred initial temperature of the process is about −55° C. to about 60° C. In general, the lower limit of the operating temperature of the reaction process is the melting point of MCS in the solvent and the upper temperature limit is determined by engineering conditions such as to avoid product decomposition and reduced efficiency of the process. In the case of anisole, depending on how much MCS is added, there is a considerable melting point depression below the anisole melting point of −37.3° C. The melting point of a given concentration of MCS in a particular solvent is easily determined by one skilled in the art without undue experimentation. [0050] The solvent of the present invention acts as a heat transport medium and as a medium for dispersing ammonium chloride formed during the formation of TSA. The solvent must have all of the following characteristics: Does not react with the starting materials, intermediates or final product Has a boiling point or vapor pressure that allows for optimum distillation/product recovery. [0053] The ratio of the vapor pressure of solvent to the vapor pressure of TSA at a given temperature is about 1:5, preferably about 1:10 or less to facilitate vacuum stripping of the reaction products from the solvent. In this description, a ratio of vapor pressure of 1:10 will be considered less than a vapor pressure ratio of 1:5. Conversely, a vapor pressure ratio of 100:1 will be considered greater than a ration of 10:1. In a preferred embodiment, the solvent is anisole and at a temperature of about 20 to about 40° C. the ratio of vapor pressure for anisole to TSA is 3.5:315 which equals about 1:90. The vapor pressure ratio is an important indicator of the separation efficiency for removing TSA and DSA from the solvent by vacuum stripping or distillation. A solvent with a low vapor pressure with respect to the vapor pressure of DSA and TSA will facilitate vacuum stripping of the DSA and TSA from the reaction solvent and collecting the DSA and TSA products. [0054] A solvent with a high vapor pressure with respect to DSA and TSA will also facilitate removal of the solvent from the DSA and TSA leaving a concentrated DSA and TSA product in a storage vessel that will not collect the lower boiling higher vapor pressure solvent. The DSA and TSA products collected may then be further purified by standard techniques such as those disclosed herein and in the literature. [0055] Suitable solvents are solvents that are aprotic, non-acidic (Lewis acidic) and solvents that do not form strong hydrogen bonds (N—H a source of hydrogen bonding). Suitable solvents have a Donor Number (“DN”) between from about 6 to about 28 and preferably between from about 6 to about 24 and a solvent polarity (“ET N ”) between from about 0.1 to about 0.4. Abboud and Notario, Pure Appl. Chem., Vol. 71, No. 4, pp. 645-718, 1999 provide the definitions of DN (page 684) and ET N and ET30 (pages 672-673) as well as how to determine these values for solvents and a comprehensive listing of these values for many solvents. The entire contents of Abboud and Notario is incorporated herein by reference. Portions of Abboud and Notario further defining DN and E T N are reproduced herein under the heading Definitions and Units of Model Dependent Scales. [0056] In chemistry a Donor Number or “DN” is a quantitative measure of Lewis basicity. A donor number is defined as the negative enthalpy value for the 1:1 adduct formation between a Lewis base and the standard Lewis acid SbCl 5 (antimony pentachloride), in dilute solution in the noncoordinating solvent 1,2-dichloroethane with a zero DN. The units are kilocalories per mole. The donor number is a measure of the ability of a solvent to solvate cations and Lewis acids. The method was developed by V. Gutmann in 1976, “Solvent effects on the reactivities of organometallic compounds”. Coord. Chem. Rev. 18 (2): 225. Likewise Lewis acids are characterized by acceptor numbers. In summary, the Gutmann Acceptor (AN) and Donor Number (DN) are measures of the strength of solvents as Lewis acids or bases. The Acceptor Number is based on the 31 P-NMR chemical shift of triethylphosphine oxide in the solvent. The Donor Number is based on the heat of reaction between the ‘solvent’ and SbCl 5 in CH 2 ClCH 2 Cl. [0057] It is believed that the increased efficiency of the synthesis of trisilylamine in the inventive solvents may be a result of the following solvent properties. An electron donor group places electron density on the silicon atom. Since chlorine is electronegative and is assigned a partial negative charge, additional electron density on the silicon atom will lengthen the S1-C1 bond, making it weaker and easier to break. If the bond is easier to break, the reaction proceeds more readily and can outcompete side reaction that lead to undesired by-products. [0058] Preferred solvents are selected from the group consisting of aliphatic hydrocarbons, aromatic hydrocarbons, symmetric ethers, unsymmetric ethers, poly-ethers, crown ethers and chloro-fluorocarbons, high-boiling point, mono-oxygenated ethers non limiting examples are; R—O—R′; R=R′; and R≠R′, wherein R and R′ are linear, branched or cyclic alkyl groups. Mixtures of the solvents are suitable in the present inventive process. The boiling point, and therefore vapor pressure of preferred solvents would be either high or low relative to TSA. In examples, solvents were selected that had at least about a 1:10 vapor pressure ratio (solvent:TSA) relative to TSA in which TSA could easily be vacuum-stripped with little solvent transport. For example, TSA has a vapor pressure of 315 ton at 25° C., whereas anisole has a vapor pressure of 3.5 at the same temperature. Donor Number (“DN”) and Solvent Polarity (“E T N ”) of Solvents [0059] [0000] Solvent type solvents DN ET30 ET N Ethers Anisole 9 37 0.198 Diethyl ether 19.2 34.5 0.117 THF 20.5 37.4 0.207 Diglyme ~20 37 0.198 Range  6-24 33-44 0.1-0.4 Aromatic Toluene 0.1 33.9 0.099 hydrocarbons Range 0.1-5   33-35 0.07-0.14 Amines Range 21-60 32-42 0.04-0.36 [0060] Several solvents were used in the examples. A preferred solvent is anisole. A non-limiting list of solvents useful in the present invention would include: anisole (methoxybenzene), high boiling ethers; di-n-butyl ether, di-t-butyl ether, di-sec-butyl ether, di-n-hexyl ether, dioxane (two oxygens, cyclic ether), diglyme. See above and Table 2, high volatility ethers may work as well such as diethyl ether and tetrahydrofuran (“THF”); these latter ethers may be more difficult to separate from TSA due to their proximities in boiling points and vapor pressures. The high-boiling ethers are more preferred, aliphatic hydrocarbons, cyclic and ring-fused hydrocarbons, aromatic hydrocarbons, and fused aromatic compounds that have melting points below 0° C. are preferred. Mixtures of solvents are also within the scope of the invention. Example 1 Synthesis Using a Schlenk Tube Reactor FIG. 1 [0061] A 250 mL Schlenk tube 19 , fitted with an internal thermocouple probe (⅛″ o.d. stainless steel, T-type) 15 , ⅛″ o.d. stainless steel ammonia sparge tube 18 , and ¼″ o.d. HDPE tubing was charged with 100 mL of anhydrous anisole under nitrogen. The tube was placed in a temperature controlled bath. The end of the sparge tube was raised above the liquid level and the solvent was cooled to −35° C. (freezing point of anisole is −37° C.). The head space nitrogen was then removed in vacuo from the Schlenk tube, with agitation of the solvent with magnetic stir bar 17 , to a final pressure of less than 1 ton. A 7.8 L (internal volume) carbon steel cylinder, which contained 900 ton pressure MCS (26.4 g, 397 mmoles) 25 was then added (through tube 27 ) into the adjacent U-trap (not shown) with the reaction tube closed. The valve adapter (see FIG. 1 ) 29 on the Schlenk tube was then opened and the MCS in the U-trap was allowed to warm to ambient temperature, upon which it condensed into the reaction tube 19 . The tube was cooled further to −60° C. and the internal pressure dropped to approximately 63 ton. (Further cooling of the solution to −65° C. resulted in the solvent freezing.) The reaction tube was then allowed to warm to −45° C. and a stream of house nitrogen was added to clear the sparge tube 18 of any MCS for several minutes. (The internal pressure increased during this time to 510 ton.) The ammonia cylinder (440 cc sslb containing 6.7 g NH 3 , 393 mmoles; internal pressure approximately 100 psig) 10 was opened up and pressurized to valve 12 . The inert gas purge was reduced by adjusting the rotameter ( 14 , Cole-Parmer 65-mm correlated flow meter, Aluminum with SS float; PN: EW-32044-06) to a lower setting (approximately 80% flow reduction). The anhydrous ammonia feed was then started by closing 24 and opening 12 ; the ammonia pressure and flow rate were adjusted by manipulating 11 and 14 (FM setting at 50). The sparge tube was quickly submerged into the MCS/anisole solution 16 and a white precipitate was immediately formed. [0062] The entire addition process was uneventful except for the formation of a mass of NH 4 Cl salt at the mid point of the solvent level in which no vortex was observed. Some gas breakthrough was observed toward the end of the addition. Very little salt formation was observed in the headspace and virtually no salt was observed at the top of the Claisen adapter 20 . The Schlenk tube was periodically removed from the bath (not shown) and shaken to break up the ammonium chloride in the reaction zone. The ammonia feed rate was reduced (FM setting at 10) whenever this was done. The volatile components (i.e. hydrides) were vacuum stripped under dynamic vacuum (Schlenk tube at about −30° C. to about −10° C. during the course of hydride removal) through two U-traps, not shown and cooled to about −30° C. (solvent trap) and about −196° C. (TSA, MCS, silane trap) 20 minutes after the ammonia flow was shut off. A total of 16.74 g of hydride mix was collected in the −196° C. trap and less than 5 mL of solvent in the other trap. The contents of the former trap were condensed into a 440 cc stainless steel lecture bottle (“SSLB” not shown) and stored in a freezer (temperature at approximately −23° C.) until it was purified later via fractional condensation. A significant amount of ammonium chloride was left behind in the trap once the hydride mixture was removed; albeit most of the salt was left in the Schlenk reaction tube and the amount of salt in the trap was less than about a gram. Later purification, via fractional condensation into two traps cooled to −78° C. and −196° C., revealed that 170 mg of residue remained in the SSLB. A total of 6.8 g TSA (63.4 mmoles) was retained in the former trap and the latter trap contents were transferred back into the SSLB. The yield of TSA, based on the amount of ammonia “consumed” in the reaction, is 74.4% (98.8% purity checked by GC-MS analysis) with a total hydride recovery of 90+% (based on silicon content). No evidence of solvent fragmentation contamination was observed in the analysis of purified TSA. Example 2 Synthesis Using a 600 cc Parr Reactor (FIG. 2 ) [0063] A 600 cc Parr 100 reactor was charged with 200 mL of anhydrous anisole 106 under nitrogen. (The reaction apparatus is shown in FIG. 2 ). The reactor was then cooled in an ice bath (not shown) and the nitrogen removed in vacuo. Monochlorosilane (65.7 g, 987 mmol., 200 mole % excess) was charged into the reactor through a dip tube 101 . The internal pressure of the reactor was approximately 900 Torr at 0° C. The dip tube was then purged with nitrogen delivered through flow meter 111 and tube 107 to clear the line and dip tube. Anhydrous ammonia was immediately added to the reactor through the dip tube. The reactor was stirred with a stirring rod 102 throughout the entire reagent loading and reaction time at a rate of 250 rpm. The temperature and pressure was monitored via an internal K-type thermocouple 103 and a 0-60 psig pressure gauge 104 . Anhydrous ammonia (7.5 g, 440 mmol.) was added to the reactor at a rate of 140 mg/min. over the course of 54 minutes from the ammonia cylinder (440 cc sslb containing 6.7 g NH 3 , 393 mmoles; internal pressure approximately 100 psig) 110 through flow meter 109 and valve 108 . The reaction mixture was stirred at 0° C. for an additional 45 minutes and the volatiles removed under dynamic vacuum. The product gas was collected in a U-trap (not shown) held at −196° C. downstream from a solvent trap (U-trap) cooled to −35° C. Less than 2 mL of solvent was collected in the solvent trap. The product mixture was transferred to a 440 cc stainless steel lecture bottle and the contents were purified via fractional condensation using two U-traps cooled to −78 and −196° C. The contents of the −78° trap contained 9.84 g TSA (92 mmol., 83% yield) and the −196° C. trap contained excess MCS and trace silane. [0064] Examples 3 to 7 were prepared by the procedure of Example 2 under the conditions described in Table 1. The yield of each example is reported in Table 1. [0065] The first run was performed in glass, the rest in a 600 cc stirred Parr reactor. The gray boxes indicate product yields that may have contained a significant amount of solvent (toluene). [0000] TABLE 1 Summary of examples 1-7. [0066] In examples 4 and 5 below, the yield and mole percent hydrides recovered results may contain residual solvent contamination. The “% salt downstream” indicates a weight percentage of ammonium chloride that is collected in the cryo-trap from the maximum amount calculated (theoretical amount) for each experiment. [0067] The reactants are contacted in a manner that optimizes reaction conditions thereby avoiding excessive reaction conditions such as heat build up from the exothermic reaction which can result in product decomposition and the formation of synthesis byproducts, notably silane and silazane polymers. For example, the process causes the ammonium halide by product of the reaction to stay in the reactor while the gaseous products such as disilylamine and trisilylamine are vacuum stripped from the solvent mixture and flow out of the reactor and are collected in a cold trap vessel substantially free of ammonium halide and solvent which can cause decomposition of the hydride products. The ammonium halide byproduct of the synthesis is crystalline under reaction conditions, therefore it remains in the solvent in the reactor while the gaseous products continue to travel up the reactor and out of the reactor. The boiling point of trisilylamine is 52° C. at one atmosphere. [0068] The reactor is run at reduced pressure or at pressures up to about 2000 Torr. Preferably the reactor is kept at pressure of about equal to or lower than the vapor pressure of the monohalosilane at any given reaction temperature. In operation, the reactor pressure will drop as the monohalosilane is depleted. Preferably maintaining the reactor internal pressure at about 100 torr to about 1500 ton. A preferred operating pressure would be about two atmospheres or less. Maximum operating pressure is about 80 psig. [0069] The present invention is directed to a process for preparing trisilylamine comprising: (a) adding a solvent to a reactor vessel; (b) adding monochlorosilane into the solvent to form a solution; (c) adding anhydrous ammonia into the solution to form a reaction mixture; (d) forming trisilylamine in the reaction mixture; (e) separating the trisilylamine from the reaction mixture; and (f) purifying the trisilylamine [0076] wherein the solvent has a DN between about 6 to about 24 and an E T N from about 0.1 to about 0.4. [0077] After the solvent is added to the reaction vessel, the temperature of the solvent may optionally be adjusted prior to condensing monochlorosilane into the solvent to form a solution. The temperature of the solvent may be adjusted to be between from about 70° C. to about −78° C., preferably from about 60° C. to about −20° C., and most preferably from about 50° C. to about −20° C. [0078] Monohalosilanes useful in the present invention include monofluorosilane, monochlorosilane, monobromosilane and monoiodosilane. Monochlorosilane is preferred. Summary Examples 8 to 13 [0079] TSA was synthesized in a 4 L Autoclave stirred-tank reactor with anisole as the solvent media. A total of six runs were conducted with varying target reaction temperatures, excess MCS amounts and the solvent to NH 3 ratio. Based on the results of the runs the following reaction conditions are recommended: [0080] Reaction Temperature equals about 20° C. to about 60° C. [0081] Excess MCS amount equals about 25% to about 40% excess to theoretical MCS amount on a mole to mole basis. [0082] TSA results are reported as a percent of theoretical yield. [0083] Solvent to NH 3 mass ratio equals about 25:1 to about 30:1. Solvent to NH 3 mass ratio will be expressed as a whole number throughout this specification. [0084] In the six runs, vacuum stripping was done from the reactor (typically at 10 to 18 psia) to a receiver in a liquid nitrogen dewar. The stripping rate was about 2.2 gm/min through a ¼ inch line and standard cylinder valve opening. The crude product was collected in the receiver through a filter to remove any salt carry over from the reactor. Also, about 6% (by mass) if the collected crude is estimated to be carried-over solvent, salt and heavies. In FIGS. 3-7 , The X axis is time in minutes and the Y axis is temperature in ° C. for the top broken line representing temperature and the Y axis is pressure (psig) for the bottom solid line of each figure. Details on the Examples [0085] In total, six runs were conducted in the 4 L reactor. Example 8 [0086] The target reaction temperature 0° C. [0087] Excess MCS equals about 63% [0088] Solvent to NH3 mass ratio is 30 (30:1) [0089] The temperature and pressure profile in the reactor as a function of time is shown in FIG. 3 . The top broken line in FIG. 3 represents temperature (° C.) and the bottom solid line represents pressure (psig) [0090] The fluctuations seen in the temperature were attributed to poor mixing. On further analysis it was determined that the viscosity of anisole is about 33% higher at 0° C. than at 20° C. So it may be that the higher viscosity of anisole combined with the increasing amount of salt in the reactor may have contributed to the temperature fluctuations. The TSA yield was 84%. Example 9 [0091] The target reaction temperature—25° C./room temperature [0092] Excess MCS—26% [0093] Solvent to NH 3 Mass ratio=28.4 [0094] The temperature and pressure profile as a function of time is shown in FIG. 4 . [0095] The top line in FIG. 4 represents temperature (° C.) and the bottom line represents pressure (psig). The TSA yield was 85.4%. Example 10 [0096] Twice the target amount of ammonia was added. The results of this run indicated that with excess NH 3 , no TSA or MCS were produced and captured in the product receiver and that only SiH 4 and NH 3 were seen in the liquid and vapor phase. The TSA yield was 0%. [0000] Vapor Phase Liquid Phase SiH 4 83.45% 11.74% NH 3 16.55% 88.02% [0097] The initial pressure in the receiver, once warmed up, was more than 180 psig, above the vapor pressure of NH 3 at room temperature and so most of the NH 3 was seen in the liquid phase this observation indicates that (i) NH 3 and TSA react in the condensed phase to form silane and (ii) no such reaction happens in the vapor phase. Example 11 [0098] Target reaction temperature=25° C./room temperature [0099] Excess MCS=39% [0100] Solvent to NH 3 ratio=25 [0101] The pressure and temperature profile during the run is given in FIG. 5 . The top line in FIG. 5 represents temperature (° C.) and the bottom line represents pressure (psig). The TSA yield was 94.3%. [0102] The reactor vapor phase was analyzed at different times during the run and the vapor phase concentration profile is given below. [0103] The concentration of MCS decreased gradually during the run with the corresponding increase in concentration of other species such as SiH 4 , TSA and DSA. Calculation of the partial pressures of different species (see table below) indicated that the SiH 4 in the vapor phase, at least initally, is from the SiH 4 in the MCS feed. Partial Pressure as a Function of Time [0104] [0000] Time MCS DSA (min) SiH 4 (psia) (psia) (psia) TSA (psia) 18 1.79 22.73 0.00 0.00 59 2.44 16.46 0.62 0.79 110 3.66 9.82 1.15 2.66 [0105] The above results indicate that as the MCS in the liquid phase is consumed, the reaction shifts to the vapor phase with the formation of TSA increasing with time, with a corresponding decrease of MCS. The increase of SiH 4 amount can either be (i) SiH 4 in the MCS feed or (ii) decomposition of TSA in anisole due to the presence of salts. Example 12 [0106] Target reaction Temp=25° C. or room temperature [0107] Excess MCS=42% [0108] Solvent/NH 3 Mass ratio=25 [0109] NH 3 addition rate=0.5 grams/min [0110] The pressure and temperature profile during the run is as follows in FIG. 6 . The top line in FIG. 6 represents temperature (° C.) and the bottom line represents pressure (psig). The TSA Yield was 81.9. % [0111] The reactor pressure was steady at about 5 psig for most of the run but after about 120 minutes the pressure increased rapidly. Samples of vapor phase in the reactor were taken at different times during the run. [0112] The amount of silane in the vapor phase at t=0 (t=time) should be from the silane in the MCS feed. An analysis of the MCS feed showed that it contained about 1% silane and so based on the amount of MCS added it can be estimated that 1.08 gms of SiH 4 was added in the feed. An overall mass balance of silane showed that about 50% of the silane in the feed MCS is in the vapor phase and so the remaining should be solubilized in anisole. An independet set of tests conducted with MCS and anisole showed that about 66% of the SiH 4 in the MCS feed can be accounted for in the vapor phase. [0113] Partial pressure of the different species as a function of reaction time was calculated (see table below). Partial Pressure as a Function of Time [0114] [0000] DSA TSA SiH 4 amt Time (min) SiH 4 (psia) MCS (psia) (psia) (psia) (grams) 0 2.89 17.11 0.06 0.03 0.52 5 3.73 16.64 0.23 0.06 0.68 38 4.47 11.13 1.29 0.74 0.81 75 5.19 10.59 0.40 1.98 0.94 123 15.00 0.05 4.75 8.05 2.72 [0115] If the initial partial pressure of silane is subtracted from the partial pressures at different times, the differential partial pressure increases as a function of time. This is true even with a correction for an increase in reaction temperature as time progresses. This indicates that (i) there is some decomposition of TSA as more salt is formed or (ii) the silane dissolved in the solvent is slowly desorbing as the reaction proceeds. Given that at the end of the NH 3 addition, the amount of silane in the vapor phase exceeded that amount added via the MCS feed, the decomposition of TSA in the presence of salt with SiH 4 evolving has been demonstrated. Example 13 [0116] Target reaction Temp=25° C. or room temperature [0117] Excess MCS=27% [0118] Slovent/NH 3 Mass ratio=26 [0119] NH 3 addition rate=0.5 grams/min [0120] The pressure and temperature profile during the run is as follows in FIG. 6 . The top line in FIG. 7 represents temperature (° C.) and the bottom line represents pressure (psig). The TSA Yield was 50.9%. [0121] Partial pressure of different species as a function of time is given in the following table. Partial Pressure as a Function of Time [0122] [0000] DSA Time (min) SiH 4 (psia) MCS (psia) (psia) TSA (psia) NH 3 (psia) 0 1.84 17.49 0.03 0.03 0 5 1.94 16.48 0.02 0.03 0 36 2.18 15.08 0.16 0.17 0 73 2.80 13.66 0.17 0.88 0 153 15.88 0.06 5.98 12.89 3.19 183 12.20 0.04 2.56 3.19 8.11 [0123] NH 3 addition was stopped at t=153 minutes and a sample was taken at that time. The reactor contents were continuously stirred for an additional 30 minutes and a sample was taken at t=183 min. Then only the reactor contents were vacuum stripped. [0124] Again, subtracting the initial partial pressure of silane from the silane partial pressure at different times shows that the silane amount increases in the vapor phase over time, even after correcting for the temperature increase. Also, NH 3 break through was seen in this run whereas no NH 3 peaks were seen in the previous runs (Examples 9, 11 and 12). The major difference is that both the excess MCS of 25% and the solvent/NH 3 mass ratio of 25 are at the low end of the operating conditions. [0125] TSA in was analyzed by a gas chromatographic procedure. The analysis conditions are indicated below. [0000] Product to be analyzed: TSA Impurities to be analyzed: purity Carrier Helium Column/Mesh RTX-1 Length 105 m, 0.53 mm i.d. Sampling condition Gas phase 100 Torr static Liquid phase—see NOTE 1 100 Torr static Reference—see NOTE 2 Varies Oven Conditions: Initial Temp (° C.)/Time (min)  35/4.5 Ramp (° C./min) 25 Final Temp. (° C.)/Time (min) 100/3.0 Standby (overnight) Temp. (° C.) 35 Detector TCD Flow A (For TCD) 20 +/− 2 ml./min Flow B (For Ref. Gas) 20 +/− 2 ml/min Approximate Retention Times: SiH 4 1.9 min MCS 2.3 min DSA 3.6 min TSA 4.7 min Si 4 NH 11 (Other Silylamine) 7.2 min Si 5 N 2 H 14 (Other Silylamine) 7.5 min [0126] Abboud and Notario (“Abboud”) provide extensive tables listing the properties of a number of solvents. Abboud Table 1a assigns numbers, Abboud Number, to a listing of compounds. Abboud Table 2a provides the ET N for the compounds and Abboud Table 2d provides the DN for the compounds. [0127] Table 2 below compiles a non-limiting list of solvents, inventive and non-inventive. Inventive compounds will have an ET N between about 0.1 and about 0.4 and a DN between about 6 and about 28. FIG. 8 graphically represents the inventive ranges of ET N and DN and the relationship of several inventive and non-inventive solvents. [0000] TABLE 1 Donor Number (“DN”) and Solvent Polarity (“E T N ”) of Solvents Solvent Identifier (Abboud Compound Number) ET N DN Name 40 0.099 0.1 Toluene 43 0.0704 5 p-xylene 44 0.068 10 1,2,3,4-tetramethyl benzene 56 0.296 0.1 1,2-dichloroethane 110 0.17 4 iodobenzene 112 0.46 14.6 acetonitrile 114 0.364 16.6 butanenitrile 129 0.117 19.2 Diethyl Ether 130 0.102 17.8 Di-n-propyl Ether 131 0.105 19 Di-isopropyl Ether 135 0.173 19 Dibenzyl ether 141 0.198 9 Anisole 142 0.182 8 Ethyl phenyl ether 144 0.336 16 bis(2-chloroethyl) ether 148 0.244 20 Diglyme 154 0.164 6 Furan 155 0.207 20.5 THF 156 0.179 12 2-methyl THF 159 0.17 22 tetrahydropyrane 161 0.164 14.8 1,4-dioxane 283 0.043 31.7 Triethyl Amine 287 0.179 27 N,N-dimethylaniline 313 0.145 50 Diethyl Amine [0128] Based on ET N and DN a non-limiting list of solvents suitable for the present invention includes: acetonitrile, butanenitrile, diethyl ether, di-n-propyl ether, di-isopropyl ether, dibenzyl ether, anisole, ethyl phenyl ether, bis(2-chloroethyl)ether, diglyme, furan, tetrohydrofuran (THF), 2-methyl THF, tetrahydropyrane and 1,4-dioxane. Definition and Units of Model-Dependent Scales [0129] ‘Overall solvation’ scales [Abboud and Notario, Pure Appl. Chem., Vol. 71, No. 4, pp. 672-673, 1999] [0000] E T (30) and E T N [0130] These are possibly the most widely used empirical solvent ‘polarity’ scales. According to Reichardt [1,3], the E T (30) value for a specified solvent is defined as the molar transition energy (in kcal/mol) for the long wavelength electronic transition of dye 1a,2,6-diphenyl-4-(2,4,6-triphenylpyridinio)-phenolate as a solution in this solvent at 25.0° C. and at a pressure of 0.1 MPa. E T (30) is obtained from the experimentally determined vacuum wavelength of the absorption maximum of this transition (lmax) through eqn (11): [0000] E T (30)/(kcal/mol)=28 591=(λmax/nm)  (11) [0131] The long-wavelength intramolecular charge-transfer absorption band (‘solvatochromic band’ [1]) exhibits very large hypsochromic shifts with increasing solvent ‘polarity’ (in Reichardt's sense). [0132] 1a is very sparingly soluble in solvents of low polarity. This has prompted Reichardt and co-workers to develop other indicators endowed with higher solubility in these media. The much more lipophilic penta-tert-butyl-substituted derivative 1b has been found to be quite satisfactory for the purpose of extending the E T (30) scale to these solvents (Scheme 1). The quantitative link between E T (30) and E T (1b) is given by eqn [12]: [0000] E T (1 b )/(kcal/mol)=0.9424E T (30)/(kcal/mol)+1.808  (12) [0133] with n=57; r=0.9990; u=0.17 kcal/mol [0000] [0134] Equation (12) allows the indirect estimation of the E T (30) values for low polarity solvents. [0135] ‘Primary’ E T (30) values, that is, those obtained directly from the study of the electronic absorption spectrum of 1a are generally known within 0.1 kcal/mol. Simple statistical considerations indicate that [0000] ‘secondary’ values, obtained through eqn 12 are affected by an uncertainty of ˜2×0.17=0.34 kcal/mol. [0136] The E T (30) values given in this compilation come from two main sources: (i) Reichardt's 1994 review [3] and (ii) Reichardt & Schafer's [C. Reichardt, P. Schäffer. Liebigs Ann. 1579 (1995).] 1995 paper. The latter contains new data as well as some revised values for selected hygroscopic solvents. We emphasize that the absorption spectra of Dimroth-Reichardt's dyes such as 1a and 1b are known to be extremely sensitive to traces of water and other hydrogen bond donor impurities [C. Laurence, P. Nicolet, C. Reichardt. Bull. Soc. Chim. Fr. 125 (1987).]. Data for super-critical CO2 are from very recent work by Reichardt and co-workers [R. Eberhardt, S. Lo{umlaut over ( )}bbecke, B. Neidhart, C. Reichardt. Liebigs Ann. Recueil 1195 (1997)]. [0137] E T N is a dimensionless ‘normalized’ scale, defined through eqn 13: [0000] E T N (Solvent)=[ E T (Solvent)−E T (TMS)/[ E T (Water)− E T (TMS)]  (13) [0138] wherein tetramethylsilane (TMS) and water are selected as rather extreme cases of ‘polarity’. The values given in this compilation are taken from the same sources as E T (30). [0139] Regarding the use of the term solvatochromism. According to Reichardt [C. Reichardt. Solvents and Solvent Effects in Organic Chemistry, 2nd edn. VCH Weinheim, (1990). (b) C. Reichardt. Chem. Rev. 94, 2319 (1994).]: ‘The term solvatochromism is used to describe the pronounced change in position (and sometimes intensity) of an UV-Visible absorption band, accompanying a change in the polarity of the medium. This term is widely used. [0140] DN and ΔH° BF3 [Abboud and Notario, Pure Appl. Chem., Vol. 71, No. 4, pp. 684, 1999] [0141] The scale DN (‘donicity scale’) was developed in 1966 by Gutmann [V. Gutmann, E. Wychera. Inorg. Nucl. Chem. Lett. 2, 257 (1966). (b) V. Gutmann. Coordination Chemistry in Non-Aqueous Solvents. Springer, New York (1971).] and has been extensively used since then. It precedes all other ‘basicity scales’ and has played a seminal role in solution chemistry. It is defined operationally as the negative of the standard enthalpy changes, ΔH° SbCl5 for the formation of the 1:1 adduct between antimony pentachloride and electron pair donor solvents D, both in dilute solution in 1,2-dichloroethane at 25.0° C. and 0.1 MPa, according to reaction 35: [0000] D(soln)+SbCl 5 (soln)⇄D:SbCl 5 (soln)  (35) [0000] ΔH° SbCl5 is given by eqn 36: [0000] ΔH° SbCl5 =ΔH 1 1 ΔH 2   (36) [0142] wherein ΔH 1 and ΔH 2 are the enthalpy changes under standard conditions for reactions 37 and 38, respectively: [0000] D(pure liquid)+SbCl 5 (soln)⇄D:SbCl 5 (soln)  (37) [0000] and [0000] D(pure liquid)⇄D(soln)  (38) [0143] All these data were determined by calorimetric techniques. In particular, ΔH 1 was obtained by dissolving pure D into a solution containing an excess of SbCl 5 . [0144] The DN scale has been very widely used, particularly in the field of coordination chemistry. [0145] The description of illustrative and preferred embodiments of the present invention is not intended to limit the scope of the invention. Various modifications, alternative constructions and equivalents may be employed without departing from the true spirit and scope of the appended claims.

The present invention is directed to a condensed phase batch process for synthesis of trisilylamine (TSA). An improved synthesis process that incorporates a solvent to help promote a condensed-phase reaction between ammonia gas (or liquid) and liquified monochlorosilane (MCS) in good yields is described. This process facilitates the removal of the byproduct waste with little to no reactor down time, substantial reduction of down-stream solids contamination and high-purity product from first-pass distillation.

FIELD OF THE INVENTION The present invention relates a RNA virus-derived peptide with modified side chains. In particular, the present invention relates a peptide, wherein the binding affinity for viral RNAs and the cellular uptake capability are regulated by altering the side-chain properties such as length or charges thereof. BACKGROUND OF THE INVENTION Tat(47-57) is the 11-amino acid basic region of HIV Tat protein (human immunodeficiency virus transactivator of transcription protein, residues 47-57) which interacts with the transactivator response element (TAR) RNA, and is responsible for cell penetration. The Tat-TAR interaction is crucial for HIV proliferation, whereas the cell penetration capability of Tat is important for inducing viral proliferation in infected neighboring cells and inducing apoptosis of nearby healthy immune cells. Since the six Arg residues in this short basic region of Tat are critical for both TAR RNA recognition and cell penetration, Tat(47-57) would be an attractive system to simultaneously study the effect of Arg side chain length on two distinct bioactivities. For example, U.S. Pat. No. 7,056,656 disclosed a Tat derived oligourea which competes with the Tat molecule of the TAR RNA in HIV-1 for the specificity to inhibit protein-nucleic acid interactions. U.S. Patent Application No. US2009/0047272 A1 disclosed an antiviral composition, comprising a nuclease covalently attached to a target ligand (Tb)r to the inhibit the viral function by cleaving viral nucleic acids, wherein the target ligand (Tb)r may be a membrane permeating peptide comprising Arg/Lys-rich peptide, which may be selected from a peptide of YGRKKRPQRRR (SEQ ID NO: 20, HIV TAT47-57). However, above-mentioned prior arts failed to regulate both biological activity of Tat peptide-TAR RNA recognition and cellular uptake activity, resulting in a Tat peptide having a high affinity for RNAs but probably losing its uptake ability, and vice versa. Thus, Tat(47-57) related drugs, which can effectively inhibit the replication of RNA virus (e.g. HIV virus), have not been found yet in the market. SUMMARY OF THE INVENTION Based on these defects in prior art, developing and improving is necessary. Thus, an objective of the present invention provided herein is a RNA virus-derived peptide with modified side chains, and thus the binding affinity between the peptide and viral RNAs can be regulated by altering the side-chain length or charges to modify this peptide such that the peptide can effectively bind with various RNA viruses. Another objective of the present invention provided herein is a method for using the RNA virus-derived peptide, whereby a composition for inhibiting RNA virus can be prepared, and thus the viral self-replication can be effectively inhibited by the high affinity between the RNA virus-derived peptide and viral RNA. A further objective of the present invention provided herein is a RNA virus-derived peptide with modified side chains and the use thereof, wherein the cellular uptake capability is regulated by altering the side-chain properties such as length or charges thereof such that the RNA virus-derived peptide may be used as a drug delivery carrier. For these objectives, a RNA virus-derived peptide with modified side chains is provided herein, including the formula (I): wherein the n 1 , n 2 , n 3 , n 4 , n 5 and n 6 may be a positive integer no greater than 8, respectively; X is NH or O Y, Y′ and Y″ may be any amino acid side chain group, respectively; and the n 1 , n 2 , n 3 , n 4 , n 5 and n 6 are not 3 at the same time when X is NH. In a preferred embodiment, the n 1 , n 2 , n 3 , n 4 , n 5 and n 6 may be 3 respectively when X is O. In another preferred embodiment, the RNA virus-derived peptide may be a HIV virus-derived peptide. Preferably, the HIV virus-derived peptide may be a Tat(47-57) derived peptide. In another preferred embodiment, the RNA virus-derived peptide may be Tyr-Gly-Agb-Lys-Lys-Agb-Agb-Gln-Agb-Agb-Agb-NH 2 (SEQ ID NO: 1), wherein the Agb may be (S)-2-amino-4-guanidinobutyric acid. In another preferred embodiment, the RNA virus-derived peptide may be Tyr-Gly-Agh-Lys-Lys-Agh-Agh-Gln-Agh-Agh-Agh-NH 2 (SEQ ID NO: 2), wherein the Agh may be (S)-2-amino-6-guanidinohexanoic acid. In a preferred embodiment, the n 1 , n 2 , n 3 , n 4 , n 5 and n 6 may be 2 respectively. In another preferred embodiment, the n 1 , n 2 , n 3 , n 4 , n 5 and n 6 may be 4 respectively. A method for preparing a composition for inhibiting RNA virus using a RNA virus-derived peptide with modified side chains is also provided herein, wherein the RNA virus-derived peptide may include the above-mentioned formula (I). In a preferred embodiment, the n 1 , n 2 , n 3 , n 4 , n 5 and n 6 may be an integer of 1 to 4 respectively, and the n 1 , n 2 , n 3 , n 4 , n 5 and n 6 are not 3 at the same time. The present invention further provides a method for preparing a drug delivery carrier using a RNA virus-derived peptide with modified side chains, wherein the RNA virus-derived peptide may include the above-mentioned formula (I). In a preferred embodiment, n 1 , n 2 , n 3 , n 4 , n 5 or n 6 may be an integer of 1 to 8 respectively, and the n 1 , n 2 , n 3 , n 4 , n 5 and n 6 are not 3 at the same time. A composition for inhibiting RNA virus is further provided, including: the RNA virus-derived peptide with modified side chains, and a pharmaceutically acceptable adjuvant. In a preferred, the composition may be applied individually, wherein a high binding affinity of the RNA virus-derived peptide with modified side chains may be used to inhibit viral self-replication such that RNA related diseases will be treated effectively. In another preferred embodiment, the composition may further include a drug, which may be used in combination with the RNA virus-derived peptide. A drug delivery carrier prepared by the RNA virus-derived peptide is also provided. In describing and claiming the invention, the following terminology will be used in accordance with the definitions set forth below to be considered illustrative, and shall not be restricted or limited by the foregoing detailed description. The term “peptide” is used herein to refer to a chain of two or more amino acids or amino acid analogues (including non-naturally occurring amino acids), with adjacent amino acids joined by peptide (—NHCO—) bonds. Thus, the peptides of the invention include oligopeptides, polypeptides, proteins, mimetopes and peptidomimetics. As used herein, the term “amino acid” is defined as alpha amino acids, encompassing any molecule containing both amino and carboxyl functional groups. The amino acids encompassed in the present invention may include either the L or D form of the amino acid, or a racemic mixture. Moreover, the amino acids can be naturally-occurring and non-naturally occurring amino acids. Thus, a peptide of the present invention can be made from genetically encoded amino acids, naturally occurring non-genetically encoded amino acids, or synthetic amino acids. As used herein, the term “composition” includes at least one active ingredient of the RNA virus-derived peptide, one or more pharmaceutically acceptable adjuvant(s), and optionally other therapeutical agent(s). The composition may further bind other drugs to be used in combination with the active ingredient, and thereby enhance the treatment effect. The term “adjuvant” is used herein to include any of pharmaceutically acceptable ingredient, which may be compatible with other ingredients and innocuous to the desired target, including: acidifying agents, additives, adsorbents, aerosol propellants, air displacement agents, alkalizing agents, anticaking agents, anticoagulants, antimicrobial preservatives, antioxidants, antiseptics, bases, binders, buffering agents, chelating agents, coating agents, coloring agents, desiccants, detergents, diluents, disinfectants, disintegrants, dispersing agents, dissolution enhancing agents, dyes, emollients, emulsifying agents, emulsion stabilizers, fillers, film forming agents, flavor enhancers, flavoring agents, flow enhancers, gelling agents, granulating agents, humectants, lubricants, mucoadhesives, ointment bases, ointments, oleaginous vehicles, organic bases, pastille bases, pigments, plasticizers, polishing agents, preservatives, sequestering agents, skin penetrants, solubilizing agents, solvents, stabilizing agents, suppository bases, surface active agents, surfactants, suspending agents, sweetening agents, therapeutic agents, thickening agents, tonicity agents, toxicity agents, viscosity-increasing agents, water-absorbing agents, water-miscible cosolvents, water softeners, or wetting agents or the like. As used herein, the drug delivery “carrier” refers to transport substances capable of delivering therapeutical agents or diagnostic agents to target cells. In the present invention, a complex may be formed by combining the present RNA virus-derived peptide and therapeutical agents or diagnostic agents, whereby these agents may be delivered using the cellular uptake property of the RNA virus-derived peptide. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 illustrated the overlaid bright-field and fluorescence microscopy images of Jurkat cells incubated with 7 μM Flu-AghTat (panel A), Flu-ArgTat (panel B), Flu-AgbTat (panel C) for 15 minutes at 37° C. in the presence of fetal bovine serum, washed and treated with trypsin at 37° C. for 5 minutes. FIG. 2 illustrated Flow cytometry results for cellular uptake of Flu-XaaTat peptides into Jurkat cells in the presence of fetal bovine serum at 37° C. Panel A: mean cellular fluorescence upon incubation with 7 μM peptide for 15 minutes. Panel B: mean cellular fluorescence upon incubation with various peptide concentrations (7, 30, 60, 120 μM) for 15 minutes. FIG. 3 illustrated Mean percent cellular uptake of Flu-XaaTat peptides into Jurkat cells in the presence of fetal bovine serum with 7 μM (panel A) and 30 μM (panel B) peptide for 15 minutes at 37° C. FIG. 4 illustrated that cell survival results upon exposure to 7 μM (panel A) and 30 μM (panel B) peptide AghTat, ArgTat, and AgbTat for 4 hours determined by MTT assays. FIG. 5 illustrated binding affinity for Tat-peptide and TAR RNA determined by fluorescence anisotropy assay. FIG. 6 illustrated the result of dissociation constant for Tat-peptide and TAR RNA determined by gel shift assay in the presence of poly(dI-dC). FIG. 7 illustrated fluorescence microscope images for Jurkat cells which were incubated with the 30 μM Tat derived peptide at 37° C., 5% CO 2 for 15 minutes and trypsinized for 10 minutes. (panel A: 6CF-Tat-Cit49; panel B: 6CF-Tat-Cit52; panel C: 6CF-Tat-Cit53; panel D: 6CF-Tat-Cit55; panel E: 6CF-Tat-Cit56; panel F: 6CF-Tat-Cit57). FIG. 8 illustrated the dose dependent fluorescence intensity of Jurkat cells in cellular uptake assays (panel A: the relationship between the peptide concentration and the fluorescence intensity; panel B: Mean fluorescence intensity of each peptide at 30 μM concentration; panel C: Mean fluorescence intensity of each peptide at 120 μM concentration). FIG. 9 illustrated the result of Ac-Agp57-Tat-NH 2 and 10 μg/mL poly(dI-dC) determined by gel shift assay. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENT Preferred embodiments of the present invention will be described below in more detail with reference to the accompanying drawings. The present invention may, however, be embodied in different forms and should not be constructed as limited to the embodiments set forth herein. Other objectives, advantages, and novel features of the invention will become more apparent from the following detailed description when taken in conjunction with the accompanying drawings. Example 1 Modifying all Arginine Side-Chain Length in Tat(47-57) Derived Peptides to Observe Effects on Both RNA Binding Specificity and Cellular Uptake Capability Material and Method All of the chemical reagents except those indicated otherwise were purchased from Aldrich. Diisopropylethylamine (DIEA), piperidine, trifluoroacetic acid (TFA), acetic anhydride, N-methyl-N-(trimethylsilyl)trifluoroacetamide, Tween-20, and acetic anhydride were from Acros. Guanidine hydrochloride was from Fluka. Dimethylformamide (DMF), ethyl acetate, dichloromethane (DCM) and hexane were from Mallinckrodt. Methanol and acetonitrile were from Merck. Ammonium persulfate and 1,4-Dioxane were from J. T. Baker. Glycerol, boric acid, bis-acrylamide, Tris-HCl, and tris(hydroxylmethyl)-aminomethane (Tris) were from Bioshop. Organic and high performance liquid chromatography (HPLC) solvents were from Merck Taiwan. N-9-Fluorenylmethoxycarbonyl (Fmoc)-amino acids, 1-hydroxybenzotriazole (HOBO, and O-1H-benzotriazol-1-yl-1,1,3,3-tetra-methyl uronium hexafluorophosphate (HBTU) were from Novabiochem, Fmoc-PAL-PEG-PS resin was from Applied Biosystems. Reagents and solvents were used without further purification. Analytical reverse phase (RP)-HPLC was performed on an Agilent 1200 series chromatography system using a Vydac C 18 column (4.6 mm diameter, 250 mm length). Preparative RP-HPLC was performed on a Waters Breeze chromatography system using a Vydac C 4 and C 18 columns (22 mm diameter, 250 mm length). Mass spectrometry of the peptides was performed on a matrix-assisted laser desorption ionization time-of-flight (MALDI-TOF) spectrometer (Bruker Daltonics Biflex IV) using α-cyano-4-hydroxycinnamic acid as the matrix. Determination of peptide concentration was performed on a UV-Vis spectrophotometer (Jasco V-650). Circular dichroism (CD) spectra were collected on a J815 spectrometer using 1 mm pathlength cell. The CD data was reported in mean residue molar ellipticity (deg·cm 2 ·dmol −1 ). The gel shift results were imaged using a Typhoon TRIO + gel imager with the emission wavelength set at 526 nm. Fluorescence intensity was acquired on a Varian Cary Eclipse fluorescence spectrophotometer. Cells were incubated using a CO 2 incubator (Thermo Scientific, Form a steri-cycle CO 2 incubaor). Cells were counted using a hemacytometer (Reichert Bright-Line, hemacytometer 1490). The fluorescence intensity of 6-carboxy-fluorescein labeled Tat peptides was measured on a flow cytometer (Becton Dickinson, FACS Canto™ II) and the peptide-treated Jurkat cells were imaged using a inverted fluorescence microscope (Olympus, IX71). The preparation of N,N-Bis(tert-butoxycarbonyl)-guanidine The synthesis was performed according to published procedures, as described in (a) Cheng, R. P. W., Y.-J.; Wang, W.-R.; Koyack, M. J.; Suzuki, Y.; Wu, C.-H.; Yang, P.-A.; Hsu, H.-C.; Kuo, H.-T.; Girinath, P.; Fang, C.-J. Amino Acids 2011, in press; (b) Feichtinger, K.; Sings, H. L.; Baker, T. J.; Matthews, K.; Goodman, M. J. Org. Chem. 1998, 63, 8432-8439; (c) Feichtinger, K.; Zapf, C.; Sings, H. L.; Goodman, M. J. Org. Chem. 1998, 63, 3804-3805, which is incorporated herein by reference. 1,4-Dioxane (30 mL) was added to a solution of guanidine hydrochloride (2.8727 g, 30.049 mmol) and sodium hydroxide (4.9573 g, 123.9 mmol) in water (30 mL), the mixture was cooled to 0° C. using an ice water bath. Di-tert-butyl-dicarbonate (14.5296 g, 66.1302 mmol) was then added to the reaction, and the residual was washed into the reaction with another 30 mL of 1,4-dioxane. The reaction was allowed to warm to room temperature and stirred for 3 days. The reaction mixture was then concentrated under reduced pressure to dryness. The resulting white emulsion was diluted with water (60 mL) and extracted with ethyl acetate (3×60 mL). The organic layer was then extracted with 10% citric acid (60 mL), water (60 mL), and brine (60 mL) and dried over anhydrous sodium sulfate. Finally, the dried organic solution was then concentrated under reduced pressure to obtain the desired product as a white powder (5.4 g, 69.2% yield). 1 H NMR (400 MHz, CDCl 3 ) 3.46 (s, 1H), 1.46 (s, 18H); ESI-MS calculated for C 11 H 21 N 3 O 4 [MH+]=260.3, observed [MH+]=260.1. The preparation of N,N′-Di-Boc-N″-trifluoromethanesulfonyl-guanidine The synthesis was performed according to published procedures, as described in “The preparation of N,N-Bis(tert-butoxycarbonyl)-guanidine”. A solution of N,N-bis(tert-butoxycarbonyl)-guanidine (3.0414 g, 12.035 mmole) and triethylamine (2.0 mL) in anhydrous dichloromethane (60 mL) was cooled to −68° C. using a dry ice/acetone bath under an atmosphere of nitrogen. Triflic anhydride (2.1 mL, 12.590 mmole) was added dropwise (2.1 mL/30 minutes). After half of the triflic anhydride was added, the color became light brown and the reaction mixture was allowed to warm to room temperature and stirred overnight. The reaction was washed with 2 M sodium bisulfate (20 mL) and water (20 mL), and then dried over anhydrous sodium sulfate. The dry organic solution was concentrated under reduced pressure and purified by chromatography on silica gel eluted with CH 2 Cl 2 to obtained the desired product (3.1 g, 64.8%). 1 H NMR (400 MHz, CDCl 3 ) 1.51 (s, 18H); ESI-MS calculated for C 12 H 20 F 3 N 3 O 6 S [MNa+]=414.4, observed [MNa+]=414.2. The preparation of N α -Fmoc-(S)-2-amino-N ω,ω′ -di(Boc)-6-guanidinohexanoic acid (Fmoc-Agh(Boc) 2 -OH) The synthesis was performed following published procedures, as described in Cheng, R. P. W., Y.-J.; Wang, W.-R.; Koyack, M. J.; Suzuki, Y.; Wu, C.-H.; Yang, P.-A.; Hsu, H.-C.; Kuo, H.-T.; Girinath, P.; Fang, C.-J. Amino Acids 2011, in press, which is incorporated herein by reference. Fmoc-Lys-OH (1.100 g, 2.98 mmol) was suspended in anhydrous dichloromethane (6 mL) under nitrogen. N-Methyl-N-(trimethylsilyl)trifluoroacetamide (1.2 mL, 6.60 mmol) was added, and then the reaction mixture was heated to reflux until a clear solution was formed. The solution was cooled to room temperature, and N,N′-di-Boc-N″-trifluoromethanesulfonylguanidine (1.410 g, 3.60 mmol) was added followed by triethylamine (504 μL, 3.60 mmol). The reaction mixture was stirred at room temperature, and the reaction was monitored by TLC. Upon completion, the reaction mixture was diluted with dichloromethane (6 mL) and washed with 2 M sodium bisulfate and water, and dried with sodium sulfate. The dried organic solution was then concentrated under reduced pressure and purified by flash chromatography on silica gel (CH 2 Cl 2 to 95:5 CH 2 Cl 2 /methanol) to obtain the desired product as a white powder (1.347 g, 73.8% yield). R f =0.17 (95:5 CH 2 Cl 2 /methanol); m.p. 85-88° C.; [α] 25 D =16.9 (0.0099 g mL −1 CHCl 3 ); 1 H NMR (500 MHz, CDCl 3 /TMS): δ=8.439 (s, 1H), 7.749-7.257 (m, 8H), 5.638 (d, J (H,H)=7.324 Hz, 1H), 4.505 (br s, 1H), 4.371 (d, J (H,H)=5.798 Hz, 1H), 4.204 (t, J (H,H)=6.867 Hz, 1H), 3.406 (m, 1H), 3.310 (m, 1H), 1.918 (m, 1H), 1779 (m, 1H), 1.479 (s, 9H), 1.471 (s, 9H), 1.647-1.277 ppm (m, 6H); 13 C NMR (75 MHz, CDCl 3 /TMS): δ=175.416, 163.388, 156.764, 156.604, 153.759, 144.425, 144.289, 141.830, 128.217, 127.602, 125.659, 120.477, 83.924, 80.395, 67.609, 54.064, 47.713, 41.187, 32.172, 29.008, 28.719, 28.567, 22.755 ppm; IR (liquid): ν bar=3226, 2983, 1720, 1617, 1512, 1450, 1416, 1368, 1335, 1137, 1054 cm −1 ; ESI-MS calculated for C 32 H 42 N 4 O 8 [MH + ]: 611.3075, observed: 611.1; HRMS calculated for C 32 H 42 N 4 O 8 [MH + ]: 611.3075, observed: 611.3063. Peptide Synthesis Fmoc-PAL-PEG-PS (0.05 mmol) was swollen in N,N-dimethylformamide (DMF, 5 mL) for 30 minutes. The resin was then deprotected by 20% piperidine/DMF (3×8 min) and rinsed with DMF (5×1 min). A mixture of 3 equivalents of the appropriately protected Fmoc amino acid, HOBt and HBTU was dissolved in DMF (1 mL). Diisopropylethylamine (DIEA, 8 equivalents) was then added to the solution and mixed thoroughly. The solution was then applied to the resin. The vial that contained the solution was rinsed with DMF (1 mL) and added to the reaction. The first coupling was carried out for 8 hours. The 8th to 14th residues were coupled for 1.5 hours. Other residues were coupled for 45 minutes. After each coupling, the resin was washed with DMF (5×1 min). The resin was subsequently washed with DMF (5×1 min) and methanol, and was lyophilized. Solid phase guanidinylation was performed to synthesize Agb- and Agp-containing peptides. For Agb-containing peptides, the corresponding Dab(ivDde)-containing peptide was synthesized first. The resin was treated with trityl chloride to protect the fluorescein moiety. Then the ivDde protecting group was removed by suspending the resin in 2% hydrazine in DMF (4 mL, 5×8 min) and shaking at room temperature. The resin was washed with DMF (4 mL, 5×1.5 min) and lyophilized. For Agp-containing peptides, the corresponding Dap(Mtt)-containing peptide was synthesized first. The Mtt protecting group was then removed by suspending the resin in 1% CF 3 COOH in CH 2 Cl 2 (4 mL, 15×3 min) and shaking at room temperature. Deprotection was continued until the filtrate no longer appeared yellow. The resin was washed with CH 2 Cl 2 (4 mL, 5×1.5 min) and lyophilized. After removal of orthogonal protecting groups from the resin-bound protected peptides, the resin was resuspended in a solution of N,N′-di-Boc-N″-trifluoromethanesulfonylguanidine (820.9 mg, 2 mmol) and Et 3 N (480 μL, 6.5 mmol) in CH 2 Cl 2 . The reaction was shaken at room temperature. The reaction was microwaved once every hour (3×7 sec, 30% power). Reaction was monitored by cleaving a small amount (about 5 mg) of peptide-bound resin and analyzed by RP-HPLC. Peptides were deprotected and cleaved off the resin by treating the resin with 95:5 trifluoroacetic acid (950 μL)/triisopropylsilane (50 μL) and shaken for 2 hours. The solution was then filtered through glass wool and the resin was washed with TFA (3×1 mL). The combined filtrate was evaporated by a gentle stream of N 2 . The resulting material was washed with hexanes, dissolved in water, and lyophilized. The peptide (1 mg mL −1 aqueous solution) was analyzed using analytical RP-HPLC on a C18 column with a flow rate of 1 mL min −1 , temperature 25° C., linear 1% min −1 gradient from 100% A to 0% A (solvent A: 99.9% water, 0.1% TFA; solvent B: 90% acetonitrile, 10% water, 0.1% TFA). Appropriate linear solvent A/solvent B gradients were used for purification on preparative RP-HPLC on C 4 and C 18 column. The identity of the peptide was confirmed by MALDI-TOF. The Peptide Preparation Example 1 AghTat(Ac-Tyr-Gly-Agh-Lys-Lys-Agh-Agh-Gln-Agh-Agh-Agh-NH 2 ) (i.e. n 1 , n 2 , n 3 , n 4 , n 5 and n 6 of formula (I) were 4 respectively; SEQ ID NO.: 2) The peptide was synthesized using 242.4 mg (0.05090 mmol) of Fmoc-PAL-PEG-PS resin. The synthesis gave 349.1 mg (66.6% yield). The cleavage yielded 119.4 mg of crude peptide (>99%). The peptide was purified by preparative RP-HPLC using C4 and C18 columns to 98.5% purity, using linear gradients PLG00 — 06 and PLG01 — 11, respectively. Retention time on analytical RP-HPLC was 17.0 minutes. The identity of the peptide was confirmed by MALDI-TOF mass spectrometry. Calculated for C 72 H 133 N 33 O 14 [MH] + : 1684.07; observed m/z: 1684.938. The Peptide Preparation Example 2 ArgTat(Ac-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln-Arg-Arg-Arg-NH 2 ) (a naturally-occurring Tat-derived peptide; i.e. n 1 , n 2 , n 3 , n 4 , n 5 and n 6 of formula (I) were 3, respectively; SEQ ID NO.: 3) The naturally-occurring Tat-derived peptide was obtained by capping a Tat(47-57) peptide (-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln-Arg-Arg-Arg-(SEQ ID NO.: 4)) at both termini and synthesized by the following method: The peptide was synthesized using 234.9 mg (0.0433 mmol) of Fmoc-PAL-PEG-PS resin. The synthesis gave 334.0 mg (59.6% yield). The cleavage yielded 116.8 mg of crude peptide (>99%). The peptide was purified by preparative RP-HPLC using C4 and C18 columns to 98.7% purity, using linear gradients PLG00 — 05 and PLG04 — 16, respectively. Retention time on analytical RP-HPLC was 15.4 minutes. The identity of the peptide was confirmed by MALDI-TOF mass spectrometry. Calculated for C 72 H 133 N 33 O 14 [MH] + : 1599.98; observed m/z: 1600.34. The Peptide Preparation Example 3 AgbTat(Ac-Tyr-Gly-Agb-Lys-Lys-Agb-Agb-Gln-Agb-Agb-Agb-NH 2 ) (i.e. n 1 , n 2 , n 3 , n 4 , n 5 and n 6 of formula (I) were 2, respectively; SEQ ID NO.: 1) The peptide was synthesized using 245.3 mg (0.05151 mmol) of Fmoc-PAL-PEG-PS resin. The synthesis gave 327.5 mg (53.6% yield). The cleavage yielded 87.0 mg of crude peptide (>99%). The peptide was purified by preparative RP-HPLC using a C4 and C18 columns to 99.0% purity, using linear gradients PLG00 — 05 and PLG03 — 14, respectively. Retention time on analytical RP-HPLC was 15.0 minutes. The identity of the peptide was confirmed by MALDI-TOF mass spectrometry. Calculated for C 60 H 110 N 33 O 14 [MH] + : 1516.89; observed m/z: 1516.72. The Peptide Preparation Example 4 Flu-AghTat Flu-AghTat (6-carboxyfluorescein-βAla-Tyr-Gly-Agh-Lys-Lys- Agh-Agh-Gln-Agh-Agh-Agh-NH 2 ; SEQ ID NO.: 4) The peptide was synthesized using 245.8 mg (0.0516 mmol) of Fmoc-PAL-PEG-PS resin. The synthesis gave 340.6 mg (52.0% yield). The cleavage yielded 122.1 mg of crude peptide (>99%). The peptide was purified by preparative RP-HPLC using C4 and C18 columns to 98.6% purity, using linear gradients PLG04 — 13 and PLG16 — 26, respectively. Retention time on analytical RP-HPLC was 26.2 minutes. The identity of the peptide was confirmed by MALDI-TOF mass spectrometry. Calculated for C 94 H 146 N 34 O 20 [MH] + : 2072.15; observed m/z: 2071.98. The Peptide Preparation Example 5 Flu-ArgTat(6-carboxyfluorescein-βAla-Tyr-Gly-Arg- Lys-Lys-Arg-Arg-Gln-Arg-Arg-Arg-NH 2 ; SEQ ID NO.: 5) The peptide was synthesized using 144.6 mg (0.0305 mmol) of Fmoc-PAL-PEG-PS resin. The synthesis gave 200.1 mg (48.3% yield). The cleavage yielded 47.5 mg of crude peptide (>99%). The peptide was purified by preparative RP-HPLC using C4 and C18 columns to 99.2% purity, using linear gradients PLG06 — 13 and PLG13 — 26, respectively. Retention time on analytical RP-HPLC was 25.7 minutes. The identity of the peptide was confirmed by MALDI-TOF mass spectrometry. Calculated for C 88 H 134 N 34 O 20 [MH] + : 1987.051; observed m/z: 1987.114. The Peptide Preparation Example 6 Flu-AgbTat (6-carboxyfluorescein-βAla-Tyr-Gly-Agb-Lys-Lys- Agb-Agb-Gln-Agb-Agb-Agb-NH 2 ; SEQ ID NO.: 6). The peptide was synthesized using 247.8 mg (0.05203 mmol) of Fmoc-PAL-PEG-PS resin. The synthesis gave 329.7 mg (46.5% yield). The cleavage yielded 136.2 mg of crude peptide (>99%). The peptide was purified by preparative RP-HPLC using C4 and C18 columns to 96.1% purity, using linear gradients PLG03 — 12 and PLG14 — 25, respectively. Retention time on analytical RP-HPLC was 27.5 minutes. The identity of the peptide was confirmed by MALDI-TOF mass spectrometry. Calculated for C 82 H 123 N 34 O 20 [MH] + : 1903.96; observed m/z: 1902.39. Peptide Concentration Determination by UV-Vis Spectroscopy The concentration of the XaaTat peptide stock solutions was determined using the Edelhoch method (The Xaa herein was Agh, Arg or Agb, respectively). A 10 mM solution was prepared for each XaaTat peptide based on weight. The UV data was collected using a 1 mm pathlength cell. The concentration of XaaTat peptide stock solutions was determined by the tyrosine absorbance in 6 M guanidinium chloride (ε 282 =1220, ε 280 =1285, ε 278 =1395, ε 276 =1455). The UV absorbance was collected at wavelengths 276, 278, 280, and 282 nm over 1 minute (60×1 sec) to accurately determine the concentration of the sample. The concentration of the peptide solutions was derived using Kaleidagraph 3.52 (Synergy Software, CA). A 7 mM solution was prepared for each Flu-XaaTat peptide based on weight (The Xaa herein was Agh, Arg or Agb, respectively). The UV data was collected using a 1 mm pathlength cell. The concentration of Flu-XaaTat peptide stock solutions was determined in pH 9 buffer (1 mM sodium phosphate, 1 mM sodium citrate, and 1 mM sodium borate) based on the absorbance of 6-carboxyfluorescein (ε 492 =81000). The UV absorbance was collected at 492 nm over 1 minute (60×1 sec) to accurately determine the concentration of the sample. Based on the Beer-Lambert law, the concentration of the peptide solutions was derived using Kaleidagraph 3.52 (Synergy Software, CA). Fluorescence Anisotropy The binding affinity of XaaTat peptides (Xaa=Agh, Arg and Agb) for TAR RNA was measured by fluorescence anisotropy using a Varian Cary Eclipse fluorescence spectrophotometer. The anisotropy (r) was calculated using the following equation: R = I VV - I VH × I HV I HH I VV + 2 ⁢ ( I VH × I HV I HH ) where I vv , I vh , I hh , I hv are the fluorescence intensities with the excitation polarizer and emission polarizer is oriented perpendicular (I vh and I hv ) and parallel (I vv and I hh ) to the direction of the polarized excitation. The fluorescein-labeled TAR RNA (F-TAR-RNA) was excited at 490 nm and the polarized emissions were monitored at 512 nm. The slits were set at 10 nm for both excitation and emission. The integration time was 20 seconds. All data were measured in a Sub-Micro (Starna Cell, Inc.) with a starting volume of 160 μL. Each addition (1 μL) of XaaTat was equilibrated for two minutes before the fluorescence signal was recorded. All experiments were performed at room temperature in TKT buffer (TKT: 50 mM Tris-HCl, pH 7.4, 20 mM KCl and 0.02% Tween20). The initial concentration of F-TAR-RNA was 25 nM. The TKT solution and peptide stock solution were both equilibrated at room temperature more than 20 minutes before the experiment. The experiment was repeated three times independently. For every single point, three measurements were performed and the average value was used for deriving the apparent (apparent K D ). The apparent for the peptide-RNA complex was derived by fitting the data assuming a 1:1 peptide-RNA stoichiometry. The fitting was performed using Kaleidagraph 3.52. (Synergy Software, CA). Gel Shift Assay The fluorescein-labeled TAR RNA (F-TAR-RNA) was purchased from Sigma. F-TAR-RNA was dissolved in diethyl pyrocarbonate treated H 2 O to give a 50 μM solution. Binding assays were performed at room temperature. Peptide and RNA were incubated in pH 7.4 buffer (10 μL) containing Tris-HCl (50 mM), KCl (50 mM), poly-dIdC (10 μg/mL), 2% glycerol, and Triton X-100 (0.05%). The F-TAR-RNA concentration was 100 nM. The samples were analyzed by loading into 12% native polyacrylamide gels in 0.5% TB buffer and electrophoresis was performed with 140 V at room temperature. Dried gels were scanned by the Typhoon TRIO+ Variable Mode Imager. Bands corresponding to the free and bound RNA were quantified using ImageQuant software. The apparent dissociation constants were globally derived from the quantified data assuming a 1:1 binding stoichiometry. Fluorescence Microscope of Live Jurkat Cells Live Jurkat cells (8×10 5 cells) were incubated with 7 μM Flu-XaaTat for 15 minutes at 37° C. in RPMI medium with fetal bovine serum. After incubation, the suspension was centrifuged at 2200 rpm. The cells were washed with phosphate buffered saline (137 mM NaCl, 2.7 mM KCl, 10 mM sodium phosphate dibasic, 2 mM potassium phosphate monobasic, pH 7.4) twice to remove extracellular peptide. To ensure removal of any surface-bound peptide, cells were then incubated with 0.05% trypsin at 37° C. for 5 minutes, washed with PBS twice. Finally, cells were resuspended in 500 μL PBS containing 1 μg/mL propidium iodide. The cells were then examined by a fluorescence microscope (Olympus IX71) (492/517 nm; excitation/observation). Cellular Uptake Assay All apparatuses were sterilized by autoclave sterilizer, and the experimental surfaces were wiped with 70% ethanol. All operations were performed in a laminar flow hood. The number of cells was determined by a hemacytometer. There were 9 squares with 1.0 mm 2 area and 0.1 mm depth in the hemacytometer. The number of cells in the four corner squares was counted and averaged. The average was multiplied by 10 4 to obtain the number of cells in a 1 mL suspension. The cellular uptake experiments were performed using 8×10 5 cells. Jurkat cells were incubated with the peptides at various concentrations (7, 30, 60, and 120 μM) at 37° C. with 5% CO 2 for 15 minutes. The cells were then washed with PBS (2 g/L KCl, 2 g/L KH 2 PO 4 , 80 g/L NaCl, 11.5 g/L Na 2 PO 4 ) (2×400 μL) to remove the fetal bovine serum which might interfere with the proteolytic activity of trypsin. The cells were then incubated with 0.05% trypsin/EDTA in PBS for 5 minutes to remove the peptides which adhered to the cell surface rather than entry to the cell. The cells were washed with PBS (2×400 μL). The cells were then resuspended in 500λ PBS and transferred into the flow tube. The cells were terminated by adding Triton-X 100 to give the dead control group. Propidium iodide (PI) was added to all samples to stain the dead cells but should not stain the live cells. Fluorescence analysis for the Jurkat cells was performed with a flow cytometer (FACScan, Becton Dickinson Bioscience). The voltage of the photomultiplier tube for forward scatter, side scatter and propidium iodide was set to 270, 470, and 350, respectively. Live cells containing appropriate forward scatter and side scatter values were selected and gated as the P1 region for normal and live cells in the live control group. The minimum propidium iodide fluorescence intensity for the dead cells treated with propidium iodide in the P1 region was set as the threshold value for dead cells. In other words, cells with propidium iodide fluorescence below the threshold value would be deemed live cells. The fluorescence of 6-carboxyfluorescein was considered when the cell morphology was in the P1 region and the propidium iodide fluorescence intensity was lower than the threshold value. The 6-carboxyfluorescein fluorescence intensity was acquired for 10,000 events at room temperature. The data presented are the mean fluorescence intensity for the 10,000 cells. Each experiment was independently repeated at least three times. MTT Assay for Determining Cell Survival The cells were grown to appropriate cell density (roughly 2×10 6 cells/mL) and transferred into new media one day before performing the experiment. For each assay, 8×10 5 cells as determined by a hemocytometer were added into an eppendorf. The cells were centrifuged at 2200 rpm for 5 minutes. The supernatant was removed by suction. Then 200 μL of 7 μM or 30 μM XaaTat peptide was added. The cells were incubated for 4 hours at 37° C. and then centrifuged at 2200 rpm for 5 minutes. The cells were washed with RPMI medium twice, and then 200 μL of 0.5 mg/mL MTT in serum free RPMI buffer was added and incubated at 37° C. for 3.5 hours. The cells were then centrifuged at 4000 rpm for 5 minutes. The supernatant was then transferred to two wells (100 μL each) of a 96-well plate. Then 400 μL DMSO was added to the cells in the assaying eppendorf. The eppendorf was then vortexed for 5 minutes to solubilize the MTT purple crystals, and centrifuged at 2200 rpm for 5 minutes. Then the supernatant was transferred to 4 wells (100 μL each) of the 96-well plate for analysis. The absorbance at 570 was determined for each well using a microplate reader with the absorbance at 655 nm used for background correction. A standard curve was generated using different number of cells (8×10 5 , 4×10 5 and 2×10 5 ). The absorbance for each sample for each peptide at each concentration was then correlated to the number of live cells based on the standard curve. Results 1. Results of RNA Binding Specificity The naturally-occurring Tat(47-57) was capped at both termini to give peptide ArgTat. All six Arg residues were replaced with Agh (one methylene longer than Arg) and Agb (one methylene shorter than Arg) to give peptides AghTat and AgbTat, respectively. To enable the detection of cellular uptake, the peptides were capped with fluorescein at the N-terminus. All peptides were synthesized by solid phase peptide synthesis using Fmoc-based chemistry, purified by reverse phase high performance liquid chromatography to greater than 95% purity, and confirmed by matrix assisted laser desorption ionization mass spectrometry. The effect of Arg side-chain length on RNA recognition for XaaTat (Xaa=Agh, Arg and Agb) peptides was investigated in the absence and presence of poly(dI-dC) by fluorescence anisotropy and electrophoretic mobility shift assays (EMSA) (See Table 1), respectively. The HIV TAR RNA was labeled with fluorescein at the 3′-terminus to enable these experiments. The apparent dissociation constants (K D ) were derived from the experimental data, suggesting that the smaller K D is, the stronger the binding affinity is. Table 1 showed that the binding affinity of AghTat and TAR-RNA and that of AgbTat and TAR-RNA had no apparent difference compared with the binding affinity of naturally-occurring derived ArgTat and TAR-RNA in the absence and presence of poly(dI-dC). This result suggested that RNA virus-derived peptides with modified side chains (AghTat and AgbTat) also have the excellent binding affinity for TAR-RNA. Moreover, to determine the specificity for peptides to TAR RNA, the strength of the binding affinity was determined in the presence of competing negatively charged poly(dI-dC). The binding affinity of AghTat was significantly reduced upon adding poly(dI-dC) in Table 1, whereas the affinity of ArgTat and TAR RNA was somewhat affected by the presence of poly-anionic poly(dI-dC). Surprisingly, the affinity of AgbTat was not affected by the presence of poly(dI-dC). That is, AgbTat had the better binding specificity for TAR RNA. These results suggested that altering the Arg side-chain length did not affect the affinity between Tat peptides and TAR-RNA, but affects specificity between them. TABLE 1 Apparent dissociation constants (K D ) for the binding of Xaa-Tat (Xaa = Agh, Arg and Agb) with HIV TAR RNA in the absence and presence of poly(dI-dC) Apparent K D [nM] a HIV TAR RNA with Peptide HIV TAR RNA b poly(dI-dC) c AghTat 36 ± 9 450 ± 40 ArgTat 25 ± 3  67 ± 19 AgbTat 32 ± 9 31 ± 9 a The apparent dissociation constants were derived from the experimental data assuming a 1:1 binding stoichiometry. b Values determined by fluorescence anisotropy experiments. The experiments were performed by titrating the peptide into 25 nM fluorescein-labeled HIV TAR RNA. c Values determined by electrophoretic mobility shift assays (EMSA). The assays were performed with 100 nM fluorescein-labeled HIV TAR RNA, varying amounts of peptide, in the presence of 10 μg/mL poly(dI-dC). 2. Results of Cellular Uptake Experiments Cellular uptake experiments were performed on Jurkat cells, because these cells belonged to a CD4+ helper T cell cancer cell line which was the target of HIV. Jurkat cells were incubated separately with various concentration (7, 30, 60, 120 μM) of Flu-XaaTat for 15 minutes at 37° C. in the presence of fetal bovine serum, and then treated with trypsin to remove cell-surface bound peptide. FIG. 1 illustrated the overlaid bright-field and fluorescence microscopy images of Jurkat cells incubated with 7 μM Flu-AghTat (panel A), Flu-ArgTat (panel B), Flu-AgbTat (panel C) for 15 minutes at 37° C. in the presence of fetal bovine serum, washed and treated with trypsin at 37° C. for 5 minutes. Obvious cellular uptake for all three peptides was shown in FIG. 1 . Cellular uptake was then investigated quantitatively using flow cytometry. FIG. 2 illustrated Flow cytometry results for cellular uptake of Flu-XaaTat peptides into Jurkat cells in the presence of fetal bovine serum at 37° C., wherein Panel A: mean cellular fluorescence upon incubation with 7 μM peptide for 15 minutes. Panel B: mean cellular fluorescence upon incubation with various peptide concentrations (7, 30, 60, 120 μM) for 15 minutes. Incubating with 7 μM peptide, Flu-AgbTat exhibited 3 times higher uptake into cells compared to Flu-ArgTat and Flu-AghTat ( FIG. 2A ). The better uptake of the Agb-containing peptide was also present at higher peptide concentrations ( FIG. 2B ). Meanwhile, Flu-AghTat exhibited more cellular uptake compared to Flu-ArgTat at concentrations higher than 7 μM. FIG. 3 illustrated Mean percent cellular uptake of Flu-XaaTat peptides into Jurkat cells in the presence of fetal bovine serum with 7 μM (panel A) and 30 μM (panel B) peptide for 15 minutes at 37° C. More than 70% of the cells showed uptake upon shortening the Arg side chain length by one methylene to Agb ( FIG. 3A ); this is up to four times the number of cells with uptake compared to the Arg peptide. On the other hand, nearly all cells exhibited peptide uptake upon raising the peptide concentration to 30 μM regardless of side chain length ( FIG. 3B ). FIGS. 2 and 3 demonstrated that both modified Flu-AgbTat and Flu-AghTat exhibited cellular uptake, which could be increased by changing Arg side-chain length in Tat(47-57). Especially when the Arg side-chain length was shortened by one methylene to Agb, it exhibited an excellent effect on uptake. 3. Results of Cytotoxicity Experiments FIG. 4 illustrated that cell survival results upon exposure to 7 μM (panel A) and 30 μM (panel B) peptide AghTat, ArgTat, and AgbTat for 4 hours at 37° C. as determined by MTT assays. These results in FIG. 4 showed that all AghTat, ArgTat and AgbTat had minimal cytotoxicity as the increasing peptide concentration. For example, 30 μM of the Agb-containing peptide had great RNA binding specificity and cellular uptake activity. Further, MTT assays on Jurkat cells showed minimal cytotoxicity upon exposure to 30 μM of the Agb-containing peptide for 4 hours at 37° C., but further studies were still needed for the development of anti-HIV therapeutics or drug delivery applications. Furthermore, these results demonstrated that altering the Arg side-chain length affects both RNA binding specificity and cellular uptake activity of Tat-derived peptides, and should be a useful strategy for developing molecules with bio-medical applications. Example 2 Modifying Each Arginine Side-Chain Charges in Tat(47-57) Derived Peptides (i.e. Replacing NH 2 Groups in Arg with O Atom) Respectively to Observe Effects on Both RNA Binding Specificity and Cellular Uptake Capability The Peptide Preparation Example 1 Tat-Cit49(Ac-Tyr-Gly-Cit-Lys-Lys-Arg-Arg-Gln-Arg- Arg-Arg-NH 2 ; SEQ ID NO.: 7) The corresponding Fmoc-Tat-Cit49 peptide (SEQ ID NO: 21, Fmoc-Tyr-Gly-Cit-Lys-Lys-Arg-Arg-Gln-Arg-Arg-Arg-NH 2 ) was synthesized using 0.2404 g (0.0505 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.3256 g of resin (51.1% yield). Retention time on analytical RP-HPLC was 31.51 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1780.63]. The Fmoc group on Fmoc-Tat-Cit49 was removed selectively using 20% piperidine/DMF (5 mL, 3×8 mins). Then the resin was reacted with a solution of acetic anhydride (95 μL, 20 equivalents) and DIEA (125 μL, 20 equivalents) for 2 hours. The synthesis gave 0.3205 g of resin (50.6% yield). The cleavage yielded 102.4 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG00 — 04) and C 18 (PLG02 — 14) columns to 98.0% purity (14.9 mg). Retention time on analytical RP-HPLC was 16.18 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 66 H 120 N 32 O 15 [1601.97]; observed [1601.32]. The Peptide Preparation Example 2 Tat-Cit52(Ac-Tyr-Gly-Arg-Lys-Lys-Cit-Arg-Gln-Arg- Arg-Arg-NH 2 ; SEQ ID NO.: 8) The corresponding Fmoc-Tat-Cit52 peptide (SEQ ID NO: 22, Fmoc-Tyr-Gly-Arg-Lys-Lys-Cit-Arg-Gln-Arg-Arg-Arg-NH 2 ) was synthesized using 0.2678 g (0.0500 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.3713 g of resin (62.6% yield). Retention time on analytical RP-HPLC was 31.03 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1780.96]. The Fmoc group on Fmoc-Tat-Cit52 was removed selectively using 20% piperidine/DMF (5 mL, 3×8 mins). Then the resin was reacted with a solution of acetic anhydride (95 μL, 20 equivalents) and DIEA (125 μL, 20 equivalents) for 2 hours. The synthesis gave 0.3651 g of resin (61.7% yield). The cleavage yielded 104.0 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG00 — 04) and C 18 (PLG02 — 14) columns to 98.2% purity (12.3 mg). Retention time on analytical RP-HPLC was 15.77 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 66 H 120 N 32 O 15 [1601.97]; observed [1601.64]. The Peptide Preparation Example 3 Tat-Cit53(Ac-Tyr-Gly-Arg-Lys-Lys-Arg-Cit-Gln-Arg- Arg-Arg-NH 2 ; SEQ ID NO.: 9) The corresponding Fmoc-Tat-Cit53 peptide (SEQ ID NO: 23, Fmoc-Tyr-Gly-Arg-Lys-Lys-Arg-Cit-Gln-Arg-Arg-Arg-NH 2 ) was synthesized using 0.2987 g (0.0508 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.4090 g of resin (65.7% yield). Retention time on analytical RP-HPLC was 31.28 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1780.57]. The Fmoc group on Fmoc-Tat-Cit53 was removed selectively using 20% piperidine/DMF (5 mL, 3×8 mins). Then the resin was reacted with a solution of acetic anhydride (95 μL, 20 equivalents) and DIEA (125 μL, 20 equivalents) for 2 hours. The synthesis gave 0.4026 g of resin (64.3% yield). The cleavage yielded 112.9 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG00 — 04) and C 18 (PLG02 — 14) columns to 97.3% purity (11.7 mg). Retention time on analytical RP-HPLC was 15.97 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 66 H 120 N 32 O 15 [1601.97]; observed [1601.48]. The Peptide Preparation Example 4 Tat-Cit55(Ac-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln-Cit- Arg-Arg-NH 2 ; SEQ ID NO.:10) The corresponding Fmoc-Tat-Cit55 peptide (SEQ ID NO: 24, Fmoc-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln-Cit-Arg-Arg-NH 2 ) was synthesized using 0.2918 g (0.0496 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.3950 g of resin (63.0% yield). Retention time on analytical RP-HPLC was 33.99 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1782.98]. The Fmoc group on Fmoc-Tat-Cit55 was removed selectively using 20% piperidine/DMF (5 mL, 3×8 mins). Then the resin was reacted with a solution of acetic anhydride (95 μL, 20 equivalents) and DIEA (125 μL, 20 equivalents) for 2 hours. The synthesis gave 0.3966 g of resin (65.6% yield). The cleavage yielded 152.8 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG00 — 04) and C 18 (PLG02 — 14) columns to 97.1% purity (7.6 mg). Retention time on analytical RP-HPLC was 16.07 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 66 H 120 N 32 O 15 [ 1601.97]; observed [1601.80]. The Peptide Preparation Example 5 Tat-Cit56(Ac-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln- Arg-Cit-Arg-NH 2 ; SEQ ID NO.: 11) The corresponding Fmoc-Tat-Cit56 peptide (SEQ ID NO: 25, Fmoc-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln-Arg-Cit-Arg-NH 2 ) was synthesized using 0.2983 g (0.0507 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.4069 g of resin (64.9% yield). Retention time on analytical RP-HPLC was 33.82 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1783.39]. The Fmoc group on Fmoc-Tat-Cit56 was removed selectively using 20% piperidine/DMF (5 mL, 3×8 mins). Then the resin was reacted with a solution of acetic anhydride (95 μL, 20 equivalents) and DIEA (125 μL, 20 equivalents) for 2 hours. The synthesis gave 0.4002 g of resin (64.2% yield). The cleavage yielded 156.9 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG00 — 04) and C 18 (PLG02 — 14) columns to 96.2% purity (6.7 mg). Retention time on analytical RP-HPLC was 15.96 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 66 H 120 N 32 O 15 [ 1601.97]; observed [1601.98]. The Peptide Preparation Example 6 Tat-Cit57(Ac-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln- Arg-Arg-Cit-NH 2 ; SEQ ID NO.: 12) The corresponding Fmoc-Tat-Cit57 peptide (SEQ ID NO: 26, Fmoc-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln-Arg-Arg-Cit-NH 2 ) was synthesized using 0.2929 g (0.0498 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.4103 g of resin (75.5% yield). Retention time on analytical RP-HPLC was 32.91 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1782.71]. The Fmoc group on Fmoc-Tat-Cit57 was removed selectively using 20% piperidine/DMF (5 mL, 3×8 mins). Then the resin was reacted with a solution of acetic anhydride (95 μL, 20 equivalents) and DIEA (125 μL, 20 equivalents) for 2 hours. The synthesis gave 0.4014 g of resin (65.3% yield). The cleavage yielded 101.5 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG00 — 04) and C 18 (PLG02 — 14) columns to 97.8% purity (11.8 mg). Retention time on analytical RP-HPLC was 16.16 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 66 H 120 N 32 O 15 [1601.97]; observed [1602.18]. The above-mentioned prepared peptides were further tested by Fluorescence Anisotropy and Gel Shift Assay according to the same protocols as mentioned in the above Example 1. The Synthesis of Cell Penetration Peptide The resin (Fmoc-PAL-PEG-PS) was swollen in DMF with shaking for 30 minutes. The resin was washed with DMF (5 mL, 5×1.5 min). The Fmoc group was deprotected by 20% piperidine/DMF (5 mL, 3×8 min) and the resin was rinsed with DMF (5 mL, 5×1.5 min). The mixture of the appropriate protected amino acid (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) was dissolved in 1 mL DMF and DIEA (12 equivalents), and then applied to the resin. The vial was washed with DMF (2×1 mL) and applied to the resin and shaken. The coupling time depended on the kind (β-branch) and location of amino acid. Arginine was triple coupled, 25 minutes each time. The first amino acid was coupled for 8 hours, residues 2˜7 for 75 minutes, residues 8-15 for 90 minutes. When the coupling was complete, the resin was washed with DMF and deblocked as described earlier. A mixture of 6-carboxy-fluorescein (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) was dissolved in 1 mL DMF and DIEA (12 equivalents) and applied to the resin. The resin was shaken for 3 hours and washed with DMF (5 mL, 5×1.5 min). The resin was washed with DMF (5 mL, 5×1.5 min), methanol (5 mL), and was lyophilized overnight. Peptides was deprotected and cleaved off the resin using TFA (5 mL), triisopropylsilane (250 μL), and ethyl 1,2-dithiol (250 μL), and shaken for 2 hours. The reaction was filtered through glass wool, and washed with TFA (3×3 mL). The combined filtrate was evaporated by a gentle stream of N 2 . The residual material was washed with hexanes (3×3 mL), dissolved in water and lyophilized. The peptide was analyzed by analytical HPLC equipped with a 250 mm length C 18 column using 1 mL/min flow rate, linear 1%/min gradient from 100% A to 0% A (solvent A: 99.9% water, 0.1% TFA; solvent B: 90% acetonitrile, 10% water, 0.1% TFA), and confirmed by MALDI-TOF MS. Different linear gradients (solvent A: 99.9% water, 0.1% TFA; solvent B: 90% acetonitrile, 10% water, 0.1% TFA) were chosen to purify each peptide by preparative RP-HPLC equipped with either a C 4 or C 18 column using 10 mL/min flow rate, linear 0.5%/min gradient. The Peptide Preparation Example 1 6CF-Tat-Cit49(6CF-βAla-Tyr-Gly-Cit-Lys-Lys-Arg- Arg-Gln-Arg-Arg-Arg-NH 2 ; SEQ ID NO.: 13) The corresponding Fmoc-Tat-Cit49 peptide (SEQ ID NO: 27, Fmoc-Tyr-Gly-Cit-Lys-Lys-Arg-Arg-Gln-Arg-Arg-Arg-NH 2 ) was synthesized using 0.2218 g (0.0466 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.3005 g of resin (51.1% yield). Retention time on analytical RP-HPLC was 31.51 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1780.63]. The Fmoc group on Fmoc-Tat-Cit49 was removed using 20% piperidine/DMF (5 mL, 3×8 mins). The resin was reacted with a mixture of β-alanine (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 1.5 hours. The resin was wash with DMF and deblocked as described earlier. The resin was reacted with a mixture of 6-carboxy-fluorescein (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 3 hours. The synthesis gave 0.2905 g of resin (30.5% yield). The cleavage yielded 139.7 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG02 — 12) and C 18 (PLG13 — 25) columns to 97.5% purity (7.4 mg). Retention time on analytical RP-HPLC was 27.2 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 88 H 133 N 33 O 21 [ 1989.04]; observed [1988.76]. The Peptide Preparation Example 2 6CF-Tat-Cit52(6CF-βAla-Tyr-Gly-Arg-Lys-Lys-Cit- Arg-Gln-Arg-Arg-Arg-NH 2 ; SEQ ID NO.: 14) The corresponding Fmoc-Tat-Cit52 peptide (SEQ ID NO: 28, Fmoc-Tyr-Gly-Arg-Lys-Lys-Cit-Arg-Gln-Arg-Arg-Arg-NH 2 ) was synthesized using 0.2496 g (0.0466 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.3460 g of resin (62.6% yield). Retention time on analytical RP-HPLC was 31.03 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1780.92]. The Fmoc group on Fmoc-Tat-Cit52 was removed using 20% piperidine/DMF (5 mL, 3×8 mins). The resin was reacted with a mixture of β-alanine (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 1.5 hours. The resin was wash with DMF and deblocked as described earlier. The resin was reacted with a mixture of 6-carboxy-fluorescein (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 3 hours. The synthesis gave 0.3560 g of resin (53.4% yield). The cleavage yielded 145.8 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG02 — 12) and C 18 (PLG13 — 25) columns to 95.5% purity (5.3 mg). Retention time on analytical RP-HPLC was 27.3 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 88 H 133 N 33 O 21 [ 1989.04]; observed [1988.85]. The Peptide Preparation Example 3 6CF-Tat-Cit53(6CF-βAla-Tyr-Gly-Arg-Lys-Lys-Arg- Cit-Gln-Arg-Arg-Arg-NH 2 ; SEQ ID NO.: 15) The corresponding Fmoc-Tat-Cit53 peptide (SEQ ID NO: 29, Fmoc-Tyr-Gly-Arg-Lys-Lys-Arg-Cit-Gln-Arg-Arg-Arg-NH 2 ) was synthesized using 0.2776 g (0.0472 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.3800 g of resin (65.7% yield). Retention time on analytical RP-HPLC was 31.28 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1780.57]. The Fmoc group on Fmoc-Tat-Cit53 was removed using 20% piperidine/DMF (5 mL, 3×8 mins). The resin was reacted with a mixture of β-alanine (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 1.5 hours. The resin was wash with DMF and deblocked as described earlier. The resin was reacted with a mixture of 6-carboxy-fluorescein (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 3 hours. The synthesis gave 0.3777 g of resin (46.6% yield). The cleavage yielded 76.8 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG02 — 12) and C 18 (PLG13 — 25) columns to 97.3% purity (2.8 mg). Retention time on analytical RP-HPLC was 27.5 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 88 H 133 N 33 O 21 [ 1989.04]; observed [1988.91]. The Peptide Preparation Example 4 6CF-Tat-Cit55(6CF-βAla-Tyr-Gly-Arg-Lys-Lys-Arg- Arg-Gln-Cit-Arg-Arg-NH 2 ; SEQ ID NO.: 16) The corresponding Fmoc-Tat-Cit55 peptide (SEQ ID NO: 30, Fmoc-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln-Cit-Arg-Arg-NH 2 ) was synthesized using 0.2982 g (0.0507 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.4037 g of resin (63.0% yield). Retention time on analytical RP-HPLC was 33.99 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1782.98]. The Fmoc group on Fmoc-Tat-Cit55 was removed using 20% piperidine/DMF (5 mL, 3×8 mins). The resin was reacted with a mixture of β-alanine (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 1.5 hours. The resin was wash with DMF and deblocked as described earlier. The resin was reacted with a mixture of 6-carboxy-fluorescein (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 3 hours. The synthesis gave 0.3973 g of resin (57.5% yield). The cleavage yielded 95.7 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG02 — 12) and C 18 (PLG13 — 25) columns to 95.9% purity (4.8 mg). Retention time on analytical RP-HPLC was 27.2 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 88 H 133 N 33 O 21 [ 1989.04]; observed [1988.91]. The Peptide Preparation Example 5 6CF-Tat-Cit56(6CF-βAla-Tyr-Gly-Arg-Lys-Lys-Arg- Arg-Gln-Arg-Cit-Arg-NH 2 ; SEQ ID NO.: 17) The corresponding Fmoc-Tat-Cit56 peptide (SEQ ID NO: 31, Fmoc-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln-Arg-Cit-Arg-NH 2 ) was synthesized using 0.2988 g (0.0508 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.4078 g of resin (64.9% yield). Retention time on analytical RP-HPLC was 33.82 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1783.39]. The Fmoc group on Fmoc-Tat-Cit56 was removed using 20% piperidine/DMF (5 mL, 3×8 mins). The resin was reacted with a mixture of β-alanine (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 1.5 hours. The resin was wash with DMF and deblocked as described earlier. The resin was reacted with a mixture of 6-carboxy-fluorescein (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 3 hours. The synthesis gave 0.4130 g of resin (64.2% yield). The cleavage yielded 183.1 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG02 — 12) and C 18 (PLG13 — 25) columns to 96.4% purity (7.0 mg). Retention time on analytical RP-HPLC was 27.1 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 88 H 133 N 33 O 21 [ 1989.04]; observed [1988.95]. The Peptide Preparation Example 6 6CF-Tat-Cit57(6CF-βAla-Tyr-Gly-Arg-Lys-Lys-Arg- Arg-Gln-Arg-Arg-Cit-NH 2 ; SEQ ID NO.: 18) The corresponding Fmoc-Tat-Cit57 peptide (SEQ ID NO: 32, Fmoc-Tyr-Gly-Arg-Lys-Lys-Arg-Arg-Gln-Arg-Arg-Cit-NH 2 ) was synthesized using 0.3070 g (0.0522 mmol) Fmoc-PAL-PEG-PS. The synthesis gave 0.4371 g of resin (75.5% yield). Retention time on analytical RP-HPLC was 32.91 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 79 H 128 N 32 O 16 [ 1782.02]; observed [1782.71]. The Fmoc group on Fmoc-Tat-Cit57 was removed using 20% piperidine/DMF (5 mL, 3×8 mins). The resin was reacted with a mixture of β-alanine (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 1.5 hours. The resin was wash with DMF and deblocked as described earlier. The resin was reacted with a mixture of 6-carboxy-fluorescein (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) for 3 hours. The synthesis gave 0.4247 g of resin (68.0% yield). The cleavage yielded 180.0 mg of crude peptide (>99% yield). The peptide was purified by preparative RP-HPLC using C 4 (PLG02 — 12) and C 18 (PLG13 — 25) columns to 96.4% purity (7.0 mg). Retention time on analytical RP-HPLC was 27.1 minutes. The identity of the peptide was confirmed by MALDI-TOF MS. Calculated for [MH + ] C 88 H 133 N 33 O 21 [ 1989.04]; observed [1988.88]. The above-mentioned prepared peptides were further tested by Fluorescence Microscope of Live Jurkat Cells and Cellular Uptake Assay according to the same protocols as mentioned in the above Example 1. Results 1. Results of RNA Binding Specificity Binding affinity for Tat-peptide and TAR RNA determined by fluorescence anisotropy assay was shown in FIG. 5 and Table 2 in which the smaller K D is, the stronger the binding affinity is. FIG. 5 illustrated that all RNA virus-derived peptides (Tat-Cit49 to Tat-Cit57) exhibited the binding affinity after the charge modification in the side chains. Further, Table 2 illustrated that naturally-occurring Tat(Tat-Arg) had a greatest binding affinity to TAR-RNA, the sequence of which from strong to weak was Tat(Tat-Arg)>Cit56˜Cit57>Cit49˜Cit52˜Cit53˜Cit55. Thus, these results demonstrated that binding affinity between Tat-peptide and TAR RNA could be regulated by modifying bearing charges of peptides. TABLE 2 results of binding affinity of Tat-peptide and TAR RNA determined by fluorescence anisotropy assay. name K D (nM) Tat-Arg 25 ± 3 Tat-Cit49  94 ± 13 Tat-Cit52 116 ± 17 Tat-Cit53 102 ± 11 Tat-Cit55  99 ± 14 Tat-Cit56 55 ± 9 Tat-Cit57 49 ± 6 Table 3 and FIG. 6 illustrated the result of the dissociation constant of Tat-peptide and TAR RNA determined by gel shift assay in the presence of poly(dI-dC) at 25° C. The studies of the binding affinity for TAR-RNA in FIG. 6 showed that replacing NH 2 groups in Tat(Tat-Arg) with 0 to neutralize positive charges may weaken the binding affinity for TAR-RNA. To determine the specificity, the strength of the binding affinity was determined in the presence of competing negatively charged poly(dI-dC). The sequence of the binding affinity from strong to weak was Tat(Tat-Arg)>Cit49˜Cit56˜Cit57>Cit52˜Cit53˜Cit55. The results shown in Table 3 and FIG. 6 suggested that altering charges in Tat peptides could regulate the specificity between Tat and TAR-RNA. TABLE 3 the result of dissociation constant of Tat-peptide and TAR RNA determined by gel shift assay in the presence of poly(dI-dC) at 25° C. name K D (nM) Tat-Arg  67 ± 19 Tat-Cit49 138 ± 21 Tat-Cit52 235 ± 45 Tat-Cit53 344 ± 73 Tat-Cit55 279 ± 74 Tat-Cit56 164 ± 40 Tat-Cit57 146 ± 30 2. Results of Cellular Uptake Experiments Fluorescence microscope and flow cytometry were applied to investigate the effect on cellular uptake after replacing NH 2 with 0 (i.e. to neutralize positive charges) in Tat peptides FIG. 7 illustrated fluorescence microscope images for Jurkat cells which were incubated with the 30 μM Tat derived peptide at 37° C., 5% CO 2 for 15 minutes and trypsinized for 10 minutes. (panel A: 6CF-Tat-Cit49; panel B: 6CF-Tat-Cit52; panel C: 6CF-Tat-Cit53; panel D: 6CF-Tat-Cit55; panel E: 6CF-Tat-Cit56; panel F: 6CF-Tat-Cit57). All six peptides exhibited significant cellular uptake in theses figures. In addition, FIG. 8 illustrated the dose dependent fluorescence intensity of Jurkat cells in cellular uptake assays (panel A: the relationship between the peptide concentration and the fluorescence intensity; panel B: Mean fluorescence intensity of each peptide at 30 μM concentration; panel C: Mean fluorescence intensity of each peptide at 120 μM concentration). Table 4 illustrated the data of FIG. 8 A, showing Jurkat cells in cellular uptake assays at various concentration. FIG. 8B showed that naturally-occurring Tat(Tat-Arg) had the best cellular uptake capability in all peptides. However, when the concentration was increased up to 120 μM ( FIG. 8C ), the sequence of cellular uptake capability was Cit49>Tat(Tat-Arg)˜Cit52˜Cit53˜Cit55˜Cit56>Cit 57. These results demonstrated that all modified peptides (6CF-Tat-Cit49 to 6CF-Tat-Cit57) exhibited cellular uptake capability, which may be regulated by modifying charges in Arg side chain. TABLE 4 Mean fluorescence intensity for Tat derived peptides of Jurkat cells in cellular uptake assays. Mean fluorescence intensity (a.u.) Peptide 30 μM 60 μM 90 μM 120 μM 6CF-Tat-Arg  2170 ± 1692  5258 ± 2006  7042 ± 1475  8628 ± 3585 6CF-Tat-Cit49 1791 ± 555 6452 ± 740 13284 ± 1441 15784 ± 1567 6CF-Tat-Cit52 1221 ± 116 5332 ± 612 10040 ± 726  12378 ± 972  6CF-Tat-Cit53 429 ± 19 3417 ± 288 6499 ± 109 10861 ± 1096 6CF-Tat-Cit55  793 ± 267 5750 ± 134 9090 ± 707 10152 ± 787  6CF-Tat-Cit56 429 ± 95 2305 ± 249 6094 ± 792 11617 ± 1339 6CF-Tat-Cit57 465 ± 72 1696 ± 167 4031 ± 377 4964 ± 811 Example 3 Modifying Partial Arginine Side-Chain Length in Tat(47-57) Derived Peptides to Observe Effects on RNA Binding Affinity The Peptide Preparation Example (i.e. n 6  of formula (I) was 1; SEQ ID NO: 19) Ac-Agp57-Tat(47-57)-NH 2 (Ac-Tyr-Gly-Arg-Lys- Lys-Arg-Arg-Gln-Arg-Arg-Agp-NH 2 ) This peptide was prepared by replacing Arg in the 57th position of the naturally-occurring Tat-derived peptide (SEQ ID NO.:3) with Agp ((S)-2-amino-3-guanidinopropionic acid). The preparing method was as follows: Fmoc-PAL-PEG-PS (0.05 mmol) was swollen in N, N-dimethylformamide (DMF) for 30 minutes. The resin was washed with DMF (5 mL, 5×1.5 min). The resin was then deprotected by 20% piperidine/DMF (5 mL, 3×8 min) and rinsed with DMF (5 mL, 5×1.5 min). The mixture of the appropriate protected amino acid (3 equivalents), HOBT (3 equivalents), and HBTU (3 equivalents) was dissolved in 1 mL DMF and DIEA (12 equivalents), and then applied to the resin. The vial was washed with DMF (2×1 mL) and applied to the resin and shaken. The coupling time depended on the kind (β-branch) and location of amino acid. Arginine was triple coupled for 25 minutes each time. The first amino acid was coupled for 8 hours, residues 2˜7 for 75 minutes, residues 8˜15 for 90 minutes. When the coupling was complete, the resin was washed with DMF and deblocked as described earlier. Then a mixture of acetic anhydride (95 μL, 20 equivalents) and DIEA (125 μL, 20 equivalents) was added, followed by washing the vial which was used for containing the resin with DMF (3×1 mL). The reactants were shaken for 2 hours, washed five times with DMF (5 mL, 5×1.5 min) and methanol (5 mL), and then lyophilized overnight. For Agp-containing peptides, the corresponding Dap(Mtt)-containing peptide (Ac-Dap57-Tat(47-57)-NH 2 ) was synthesized first. The Mtt protecting group was then removed by suspending the resin in 1% CF 3 COOH in CH 2 Cl 2 (4 mL, 15×3 min) and shaking at room temperature. Deprotection was continued until the filtrate no longer appeared yellow. The resin was washed with CH 2 Cl 2 (4 mL, 5×1.5 min) and lyophilized. After removal of orthogonal protecting groups from the resin-bound protected peptides, the resin was resuspended in a solution of N,N′-di-Boc-N″-trifluoromethanesulfonylguanidine (820.9 mg, 2 mmol) and Et 3 N (480 μL, 6.5 mmol) in CH 2 Cl 2 . The reaction was shaken at room temperature. The reaction was microwaved once every hour (3×7 sec, 30% power). Reaction was monitored by cleaving a small amount (about 5 mg) of peptide-bound resin and analyzed by RP-HPLC. Peptides were deprotected and cleaved off the resin by treating the resin with trifluoroacetic acid (5 mL)/triisopropylsilane (250 μL) and shaken for 2 hours. The solution was then filtered through glass wool and the resin was washed with TFA (3×3 mL). The combined filtrate was evaporated by a gentle stream of N 2 . The resulting material was washed with hexanes (3×3 mL), dissolved in water, and lyophilized. The peptide was analyzed using analytical RP-HPLC and confirmed by MALDI-TOF MS. The analytical condition of RP-HPLC was on a 250 mm C 18 column with a flow rate of 1 mL min −1 , temperature 25° C., linear 1% min −1 gradient from 100% A to 0% A (solvent A: 99.9% water, 0.1% TFA; solvent B: 90% acetonitrile, 10% water, 0.1% TFA). Different linear gradients (solvent A: 99.9% water, 0.1% TFA; solvent B: 90% acetonitrile, 10% water, 0.1% TFA) were chosen to purify each peptide by preparative RP-HPLC equipped with either a C 4 or C 18 column using 10 mL/min flow rate, linear 0.5%/min gradient. Preparative RP-HPLC was performed on a Waters Breeze chromatography system using a Vydac C 4 and C 18 columns (22 mm diameter, 250 mm length). The dissociation constant for Tat-peptide (Ac-Agp57-Tat-NH 2 ) and TAR RNA determined by gel shift assay (The detail method referred to Example 1 and needn't have been given herein) in the presence of 10 μg/mL poly(dI-dC). FIG. 9 illustrated the result of Ac-Agp57-Tat-NH 2 and 10 μg/mL poly(dI-dC) determined by gel shift assay. From FIG. 9 , the determined K D of Ac-Agp57-Tat-NH 2 and TAR-RNA was 56±12 had no apparent difference compared with that of naturally-occurring Tat (K D =67±19), suggesting that the non-naturally occurring Tat peptides which have a side-chain length of n=1 also had the binding affinity. Conclusion: Examples of the present invention showed that altering Arg side-chain length (e.g. altering shortening Arg by one methylene to Agb in Tat-derived peptides) or charges thereof (e.g. replacing NH 2 with O) resulted in improvement in both bioactivities: RNA binding and cell penetration. These results suggested that introducing similar but previously non-existing building blocks such as non-natural amino acids can alter the bioactivity landscape. Furthermore, these results demonstrated that altering the Arg side-chain length affected both RNA binding, and thereby design anti-HIV therapeutics or drug delivery applications. The above-disclosed subject matter is to be considered illustrative, and not restrictive, and the appended claims are intended to cover all such modifications, enhancements, and other embodiments, which fall within the true spirit and scope of the present invention. Thus, to the maximum extent allowed by law, the scope of the present invention is to be determined by the broadest permissible interpretation of the following claims and their equivalents, and shall not be restricted or limited by the foregoing detailed description.

The present invention provides a RNA virus-derived peptides with modified side chains, wherein the side chains of the RNA virus-derived peptide are modified by altering the side-chain length or charges thereof such that the RNA virus-derived peptide has a high binding affinity for viral RNA and an high cellular uptake capability. The present invention also provides a composition for inhibiting RNA virus wherein the RNA virus-derived peptide can effectively inhibit viral self-replication and treat related diseases by its high affinity for viral RNA. A drug delivery carrier is also provided, wherein the RNA virus-derived peptides can carry desired drugs to the intracellular target due to its cellular uptake capability and thereby enhances the drug-delivery and treating efficiency.

[0001] This U.S. patent application claims priority to and is a continuation patent application of U.S. patent application Ser. No. 13/804,115 filed on Mar. 14, 2013, which is a continuation-in-part of U.S. patent application Ser. No. 12/504,870 filed on Jul. 17, 2009, which is incorporated herein by reference in its entirety, and which claims priority to U.S. provisional patent application Ser. No. 61/090,794 filed on Aug. 21, 2008. TECHNICAL FIELD [0002] The present invention pertains to methods and systems for introducing potential new workers to the field of welding, and more particularly, to computer generated virtual environments that simulate welding processes. BACKGROUND [0003] In recent decades, welding has become a dominant process in the manufacture and construction of various products. Applications for welding are widespread and used throughout the world for the construction of ships, buildings, bridges, vehicles, and pipe lines, to name a few examples. Many welding tasks can be automated reducing the need for skilled labor. However, automated welding applications must be set up and managed by knowledgeable welders. Other welding applications aren't confined to a factory floor. Applications, including the construction of pipe lines or buildings, are welded in the field and require the mobility of an experienced welder. Accordingly, there is ongoing need for trained personnel who can adapt to the challenges of welding processes. [0004] The demand for skilled welders remains high, despite reductions in manufacturing, in many regions of the world. In the United States, the average age of the welding professional is increasing, with many individuals approaching retirement age. Over the next decade, the number of available experienced welders is expected to significantly decline as workers retire from the welding profession. Many young people entering the workforce today are choosing advanced education over skilled trades and many of those workers entering the trades are dissuaded from a career in welding despite good working conditions. Programs and organizations promoting S.T.E.M. (Science Technology Engineering Math) and S.T.E. (Science and Technology/Engineering) education are valuable in revitalizing the interest of individuals in technology related fields. BRIEF SUMMARY [0005] The embodiments of the present invention pertain to a computer program product and processor based computing system that provides processing means for executing coded instructions and input means for interacting with said processing means to create a virtual welding environment. The system establishes an objective to change a functional or operational state of a virtual article, and directs the end user to perform at least one virtual welding operation for changing its functional state. [0006] One embodiment provides a tablet-based computing device. The tablet-based computing device includes a touch-screen display and computer memory storing at least one welding software application providing a virtual welding process. The tablet-based computing device further includes processing means operable to execute coded instructions of the at least one welding software application to generate an interactive virtual welding environment and to display the interactive virtual welding environment on the touch-screen display. The tablet-based computing device also includes an input means configured to interact with the touch-screen display when manipulated by a user to direct at least a spatial orientation of a virtual welding tool in the virtual welding environment while performing a virtual welding activity corresponding to the virtual welding process. [0007] Another embodiment provides a tablet-based computing device or tablet. The tablet includes a display, wireless communication means, and computer memory storing at least one software application. The tablet also includes processing means operable to execute coded instructions of the at least one software application. The coded instructions are executed to access at least one virtual reality welding system via the wireless communication means to download user virtual welding activity information from the at least one virtual reality welding system to the tablet. The coded instructions are also executed to generate a summary of user virtual welding progress based on the user virtual welding activity information, and display the summary of user virtual welding progress on the display. [0008] A further embodiment provides a method of virtual welding. The method includes generating a dynamic virtual welding environment within a computer-based platform and displaying stereoscopic three-dimensional (3D) images of at least a portion of the dynamic virtual welding environment on a display screen of the computer-based platform. The method further includes viewing the stereoscopic three-dimensional images using 3D glasses, resulting in at least a portion of the dynamic virtual welding environment appearing to project out of the display screen into 3D space. The method also includes virtually welding a virtual weldment of the projected portion of the dynamic virtual welding environment in 3D space using a mock welding tool while viewing the stereoscopic three-dimensional images using 3D glasses. BRIEF DESCRIPTION OF THE DRAWINGS [0009] FIG. 1 is a perspective view of a simulating device and end user according to the embodiments of the subject invention; [0010] FIG. 2 is a close up perspective view of a simulating device depicting a virtual environment according to the embodiments of the subject invention; [0011] FIG. 3 is an image of a virtual environment showing a virtual article according to the embodiments of the subject invention; [0012] FIG. 4 is an image of a virtual environment showing a virtual article and user interface screen according to the embodiments of the subject invention; [0013] FIG. 5 is an image of a virtual environment showing a user interface screen according to the embodiments of the subject invention; [0014] FIG. 6 is a block diagram depicting a method of a game played on a processor based computing device executing coded instructions; [0015] FIG. 7 is a block diagram depicting a method for training welding activity; [0016] FIG. 8 illustrates an example embodiment of a tablet-based computing device (or tablet) having a stylus; [0017] FIG. 9 illustrates an example embodiment of the tablet-based computing device of FIG. 8 (or tablet) without the stylus; [0018] FIG. 10 illustrates a schematic block diagram of an example embodiment of the tablet of FIGS. 8 and 9 ; [0019] FIG. 11 illustrates an embodiment of a tablet-based computing device in communication with a virtual reality welding system via an external communication infrastructure; [0020] FIG. 12 illustrates a schematic block diagram of an example embodiment of the tablet of FIG. 11 ; [0021] FIG. 13 illustrates a portion of a dynamic virtual welding environment, generated by a computer-based platform, that is stereoscopically projected out of a display screen of the computer-based platform into 3D space as viewed by a user wearing 3D glasses, in accordance with an embodiment; and [0022] FIG. 14 is a flowchart of an example embodiment of a method of virtual welding. DETAILED DESCRIPTION [0023] Referring now to the drawings wherein the showings are for purposes of illustrating embodiments of the invention only and not for purposes of limiting the same, FIGS. 1 and 2 shows a simulator or simulating device depicted generally at 10 . The simulating device 10 generates a virtual environment 9 that simulates a three dimensional setting, which may be an industrial or commercial setting utilizing one or more manufacturing processes. The virtual environment 9 may be depicted on an imaging device 22 as viewed by an end user 11 . In particular, the simulating device 10 may depict a virtual environment 9 that facilitates interaction between the end user 11 and one or more virtual articles 16 . An input device 13 may be included that senses activity when manipulated by the end user 11 . Data from the input device 13 may be communicated to the simulating device 10 and is used to maneuver objects with the virtual environment 9 in or near real time. In one embodiment, the simulating device 10 functions to generate one or more virtual articles 16 that can be acted upon by virtual tools 26 generated in a similar manner. It follows that the virtual tools 26 may be employed in the virtual environment 9 when the input device 13 is maneuvered by the end user 11 in the real world. [0024] The simulating device 10 may generate a virtual environment 9 having virtual articles 16 that resemble components of a particular manufacturing or construction process. In one embodiment, the virtual environment 9 may comprise a welding environment 9 a depicting one or more articles for assembly together via a welding process. Accordingly, the virtual tools 26 may comprise a welder 32 and welding torch 34 . In this manner, the simulating device 10 displays virtual articles 16 being welded together by a virtual welder 32 as interactively controlled by the end user 11 . The simulating device 10 may be realized as a training platform for exposing individuals to a particular manufacturing process, or may be realized as a game played to achieve a stated objective, both of which will be discussed further in a subsequent paragraph. It is expressly noted that while the embodiments of the present invention are described in the context of a virtual welding environment 9 a and one or more welding process, persons of skill in the art will understand its application to other industrial or commercial processes. [0025] With continued reference to FIG. 1 , the simulating device 10 may be constructed of electronic hardware comprising a processor based computing device 24 operable to run, i.e. execute, a computer program product. In one embodiment, the processor based computing device 24 comprises a microcomputer in any of various configurations including but not limited to: a laptop computer, a desktop computer, work station, server or the like. Alternatively, the processor based computing device 24 may comprise a gaming system like that manufactured by Nintendo®, Microsoft® or Sony®. In this manner, the processor based computing device 24 may be a commercially available system readily available for purchase by an end user 11 . The processor based computing device 24 may include one or more logic processor based systems 25 , or logic processors 25 , such as a programmable microprocessor, although any type of logic processor 25 may be utilized in the simulating device 10 without departing from the intended scope of coverage of the embodiments of the subject invention. The processor based computing device 24 may further include support circuitry including electronic memory, such as RAM or ROM along with other peripheral support circuitry that facilitate operation of the logic processor(s) 25 . Additionally, the processor based computing device 24 may include data storage, examples of which include hard disk drives, optical storage devices and/or flash memory for the storage and retrieval of data in a manner well known in the art. Thus, the processor based computing device 24 may be programmable and operable to execute coded instructions, as also referred to as programmed algorithms, which may be a computer program product written in a high or low level programming language. It is noted that any form of programming or type of programming language may be used to code algorithms as executed by the simulating device 10 for simulating the virtual environment 9 , 9 a. [0026] The simulating device 10 and, more particularly, the processor based computing device 24 may be communicated to and used in conjunction with other similarly or dissimilarly constructed systems. Input to and output from the simulating device 10 , termed I/O, may be facilitated in this embodiment by networking hardware including wireless as well as hard wired (directly connected) devices. Communication between simulating devices 10 , or systems, may be accomplished remotely as through a network, like a wide area network (WAN) or local area network (LAN) via network hubs, repeaters, or by any means chosen with sound judgment. Communications may be established through, but are not limited to: direct connection of multiple simulating devices 10 , web-based connectivity, virtual private networks, and/or SSL (Secure Sockets Layer) encrypted communication. It is noted that the relationship between simulating devices 10 may be peer-to-peer, client-server, or any hybrid combination thereof without departing from the scope of coverage of the embodiments of the subject invention. In this manner, information may be transmitted between systems 10 as is useful for simulating or interacting with the virtual environment 9 , 9 a . In one embodiment, network communications may be used to download virtual articles 16 or virtual tools 26 for changing the game scenario. Alternatively, new environments may be downloaded for training a different manufacturing process, the details of which will be discussed further below. It is further contemplated in another embodiment that the simulating device 10 may generate a virtual environment 9 , 9 a that may be acted upon by multiple end users 11 each working from the same system or separate systems networked together. Still, any manner of communicating one or more simulating devices 10 together may be utilized without departing from the intended scope of coverage of the embodiments of the subject invention. [0027] With continued reference to FIGS. 1 and 2 , the simulating device 10 may include an imaging device 22 for displaying the virtual environment 9 , which may be a virtual welding environment 9 a . The imaging device 22 may comprise a display screen operable to display images generated by the processor based computing device 24 and the computer program product. In one embodiment, the display screen may include a computer monitor and/or television screen comprised of CRT (Cathode Ray Tube) or LCD (Liquid Crystal Display), although any type of monitor, screen, display, or projection device may be utilized as is appropriate for the embodiments of the subject invention. Information for generating the images on the display screen may be stored in memory within the processor based computing device 24 . As memory is updated or changed during execution of the computer program product, images on the display screen may be dynamically changed in real time. Still any method or means for displaying the virtual environment 9 , 9 a on the imaging device 22 may be chosen as is appropriate for use with the embodiments of the subject invention. [0028] With reference now to FIGS. 2 and 3 , the input device 13 may function to interface activity in the real world with the virtual environment 9 . In particular, the input device 13 senses the actions of the end user 11 and translates those actions into data recognizable by the simulating device 10 . The data is communicated to the logic processors 25 and may be used to interactively engage the virtual tools 26 and/or the virtual articles 16 . In one embodiment, the computer program product processes the data and makes changes to the virtual environment 9 in real time. In this manner, as the end user 11 manipulates the input device 13 in three dimensional space, objects in the virtual environment 9 , 9 a move in a corresponding manner, i.e. in direct relationship to movement of the input device 13 . For example, the end user 11 may visualize one or more virtual objects on the imaging device 22 , including a virtual representation of a real world tool. Accordingly, the end user 11 may move the input device 13 in a particular direction resulting in a corresponding movement of the virtual object. In the welding embodiment illustrated herein, the input device 13 may represent the welding torch 34 . Movement of the input device 13 therefore translates into movement of the virtual welding torch 34 in the virtual environment 9 a . The input device 13 may include switches that activate the virtual welder thereby initiating the welding process. The end user 11 may then guide the input device 13 along a trajectory that coincides with the weld joint as depicted on the imaging device 22 . [0029] Still referring to FIG. 2 , the input device 13 may include one or more sensors 37 that detect movement and/or orientation in three dimensional space. The sensors 37 may be integrated into the input device 13 and positioned at various locations for detecting different types of activity. For example, the sensor 37 , or sensors, may detect spatial orientation, i.e. the direction that an object is pointing, as mentioned above. The sensors 37 may also detect motion in a particular direction. Additionally, the sensors 37 may detect velocity and/or acceleration, which may encompass not only the magnitude of change in position or speed, but also direction. However, any type of activity related to the input device 13 may be detected by the sensors 37 without departing from the intended scope of coverage of the embodiments of the subject invention. Examples of sensors 37 may include but are not limited to: inertial sensors like accelerometers, proximity sensors, infrared sensors, photoelectric and optical sensors, and the like. It is noted here that the input device 13 may further incorporate switching means 38 for interfacing with the virtual environment 9 . The switching means 38 may include: pushbuttons, triggers or switches. In this way, virtual activity may be initiated, interrupted, or terminated as desired by depressing or releasing the switch respectively. Illustratively, the virtual welder 34 may be turned “on” or “off” by depressing or releasing a trigger switch. It is to be construed that any type, quantity or grouping of sensors 37 or switching means 38 may be integrated into the input device 13 as chosen with sound judgment. Still other means of tracking movement of the input device 13 may be included as a separate unit that resides in an adjacent region proximal to the coupon. Position and/or orientation data generated by the tracker, i.e. tracking means, may be used in conjunction with or in place of data generated by the input device 13 . [0030] In one embodiment, the input device 13 may be commercially available for purchase and use. One example may include a manually moveable device, like a computer mouse having an optical sensor for detecting movement along an adjacent surface. Another example of input device 13 may comprise a gaming joystick or controller, which may include a connector for plugging into an I/O port or may include wireless means for communication. The Wii wireless controller manufactured by Nintendo® is one exemplary type of input device, although other commercially available controller devices may be utilized as are suitable for use with a particular processor based computing device 24 . Other embodiments contemplate customized controllers, which may be fashioned to physically look like a particular virtual tool 26 , e.g. a welding torch 34 . Interaction with the simulating device 10 is thereby enhanced by a physical object having a real world feel and look that resemble the virtual tools 26 depicted on the imaging device 22 . It is noted that the customized controller may be substantially similar in size, shape and/or weight to the real world tool for which the controller is intended to resemble. Other embodiments include an attachment that connects to the commercially available input device 13 and resembles a particular virtual tool 26 to enhance the end user's experience in interacting with the virtual environment 9 . In one embodiment, the attachment may be an overlaying component and/or a component that attaches to and extends from the input device 13 . However, it is expressly noted that any configuration of customized controller or attachment may be chosen as is appropriate for use with the embodiments of the subject invention. Accordingly, at least part of the simulating device 10 may be packaged as a kit for use with any type of processor based computing device 24 , commercially available or otherwise. In another embodiment of the subject invention, the kit may include a welding coupon that may resemble a virtual article 16 displayed in the virtual environment 9 , 9 a . Accordingly, the welding coupon may function as a guide in the real world for assisting the end user in acting in the virtual environment 9 , 9 a . The kit may also comprise tracking means like that mentioned above. In other words, a tracking unit may be provided in addition to the input device 13 for sensing the end user's 11 movement during play. [0031] With reference to FIGS. 3 , 4 and 6 , the simulating device 10 may comprise a game having a stated objective to be accomplished by the end user 11 . In one particular embodiment, the game may comprise a welding game where the objective is to weld one or more virtual articles 16 together. Achieving the objective may require the end user 11 to perform a number of welds each to a predetermined level of quality. That is to say that the game facilitates user interaction with processor based computing device 24 via the input device 13 for satisfactorily performing one or more virtual welds in the virtual welding environment 9 a . During play, the game presents the end user 11 with a scenario incorporating one or more themed virtual articles 16 a . Illustratively, the scenario may relate to motorized vehicles and may depict a number of virtual components that can be welded together to assemble a functioning motorcycle or racecar. In another exemplary scenario a jet aircraft is simulated that is in need of repairs, which may require welding before taxiing down the runway and taking off. Other examples include building structures or bridges that require repair or construction before people occupy the building or the train crosses the bridge respectively. However, any scenario theme may be chosen without departing from the intended scope of coverage of the embodiments of the subject invention. It will be realized that the game rewards successful completion of the objective, in one manner, by graphically displaying the themed article functioning in its environment, e.g. the motorcycle drives away or the train crosses the bridge. The simulating device 10 also takes in account the personal interests of the end user 11 . In one embodiment, the game gives the end user lithe option of selecting a scenario that is familiar increasing his or her level of interest. Accordingly, the game may be programmed with multiple scenario options for appeal to a broad range of individuals. [0032] From the aforementioned description, it follows that the themed virtual article 16 a of the scenario has some deficiency that requires repair or assembly before becoming operational. During game initialization, i.e. game start up, the themed virtual article 16 a may be instantiated having an inoperative state or, in other words, is created not working properly or not working at all. In the present examples, the initial “inoperative” state may be represented and simulated by one or more broken brackets, a stack of unassembled I-beams, a cracked pipe, or any repairable element fitting the scenario theme. Accomplishing the game objective therefore requires the end user 11 to interact with the virtual environment 9 a to perform virtual welding that changes the operational state of the themed virtual article 16 a . It is noted here that accomplishing the game objective may require successful completion of multiple levels of play. That is to say winning the game requires successfully changing the operational state of each virtual article 16 a in every level of play. [0033] At an introductory level, the game displays one or more virtual articles 16 that correspond to the scenario selected by the end user 11 . The end user 11 is then instructed to perform a particular type of weld relating to the deficiency of the virtual article 16 . It may be assumed that the end user 11 has little or no welding experience. Accordingly, a tutorial may be provided that presents information to the end user 11 about the welding process or welding techniques needed for achieving the objective for that level. Display of the tutorial may be initiated or controlled by the end user 11 via a graphical user interface (GUI), in one example, as selected by a “help” button. In the alternative, tutorial screens may automatically be presented if the end user's performance falls below a satisfactory level. In one exemplary manner, the instructions may be displayed in written form, an example of which may include a setup screen. Instructions may also be provided audibly and, more specifically verbally, to describe the process and/or motions needed to complete setup and a particular welding task. In either case, the instructions may be presented in one of a plurality of languages to accommodate individuals residing in different regions of the world. One embodiment is contemplated where the game graphically or pictorially presents tutorial information. In this instance, literacy of the end user 11 is not required to play the game. [0034] Game play proceeds as the end user 11 engages the input device 13 to mimic movements characteristic of performing a weld. Progression through the game may depend on how well the end user 11 performs a virtual weld, which may relate to the level of virtual weld quality. In this manner, advancing to the next level, as will be discussed further in a subsequent paragraph, requires successful completion of the previous game stage. In making that determination, one or more parameters may be measured to determine the level of virtual weld quality. In processes well known in the real world, weld quality depends on many factors like the distance between the torch tip and the weld joint, which may vary with the type of welding process, the materials being welded, the welder settings, and the like. Corresponding real world parameters may be coded into the computer program product for judging the end user's 11 performance and for determining the quality of the virtual weld. [0035] Completion of a particular game level may require the end user 11 to perform the one or more virtual welds to predetermined performance standards as determined by the computer program product. Performance parameters may be programmed into the computer program product that correlate to good welding practices and may consist of: weld torch 34 position, roll and pitch angles of orientation and travel speed. Sensor data from the input device 13 may be compared to preprogrammed parameters to determine whether or not the end user 11 has stayed within acceptable limits. In one particular embodiment, weld quality may be determined by monitoring the distance between the torch tip in relation to the center of the weld seam while maintaining proper pitch and roll angles during the virtual welding process. However, it is to be construed that other parameters may be utilized in determining if the end user 11 has successfully completed a virtual weld. [0036] In one embodiment, the simulating device 10 provides or calculates a score resulting from game play. The score, which may also be a grade, may be derived from the performance data of the end user 11 . Performance data may pertain to how well the end user 11 performs the virtual weld, that is to say how closely the end user 11 maintains the virtual tools 26 or welding torch 34 within limits for acceptable welding practices. Examples may include but should not be limited to, welding torch angle or distance to the virtual article 16 The score or grade may also be derived from end user selections made with respect to the problem-based scenarios as will be discussed further in a subsequent paragraph. [0037] Simulating device 10 may provide feedback to help the end user 11 in performing the virtual welds. In the real world, a welder receives feedback by observing the weld bead as the torch travels along the weld joint. The simulating device 10 may similarly depict a virtual weld bead correlating to the end user's movement of the virtual welding torch 34 . In one embodiment, the shape of the virtual weld bead is determined by factors including torch angle, travel speed and distance to the work piece, as well as welding power source settings, although other factors may be incorporated as is appropriate for use with the embodiments of the subject invention. In this manner, the end user 11 can make adjustments during the virtual welding process for laying down an acceptable weld bead thereby simulating real world activity. [0038] Referencing FIG. 5 , to further assist the end user 11 , performance guides 41 may be included that provide quantitative feedback on the position and orientation of the virtual welding torch 34 . In one particular embodiment, “indicating bars” 42 are included that show the pitch and roll angles of the virtual welding torch 34 . Other types of performance guides 41 are contemplated that display the distance between the torch tip and the weld joint. Additional welding parameters incorporated into the performance guide 41 will become apparent to those skilled in the art. [0039] The performance guides 41 may display the actual numerical values of the torch position, which in the current embodiment shows pitch and roll angles. The values displayed may show the angles as measured from an absolute reference like the vertical or horizontal plane. Alternatively, the performance guides 41 may display angle values that relate to offsets from the ideal torch position or orientation. The performance guide 41 may indicate which values are outside the range for achieving an acceptable weld. [0040] In one embodiment, the performance guides 41 may flash, changer color and play an audible sound that indicates when the welding torch 34 is out of position. In this way, the end user 11 , through repeated use, learns correct welding techniques. As the end user 11 gains experience, he or she will naturally maintain the welding torch 34 at the proper orientation throughout the entire welding process. At one point, it may no longer be necessary to display the performance guides 41 . Accordingly, the computer program product may be programmed to selectively turn the guides 41 “on” or “off.” [0041] As previously mentioned, the game may incorporate different levels of play. The levels may be differentiated by scenario, i.e. by changes in the themed subject articles 16 a being welded. Alternatively, the levels of play in a particular scenario may differ by the types of weld joints and/or the number of virtual article pieces to be welded together. For example, a more fundamental level may simulate welding a single lap joint embodied by overlaid frame components of building structure. Another level of play may simulate performing a pipe weld as found on a motorcycle tailpipe or pipeline. Still, other examples are contemplated wherein overhead or vertical butt joints are to be welded for repairing the frame of motor vehicle. At each game level, the welding objectives must each be performed to within predetermined quality boundaries in succession, before proceeding to the next level. In this way, basic welding skills may be taught by progressively introducing increasingly complicated weld joint configurations and more advanced welding techniques. [0042] The game objective may be accomplished when the end user 11 successfully performs, i.e. meets or exceeds predetermined limits of weld quality for, all of the virtual welds in a given scenario. That is to say that the end user 11 performs each weld on every level to a minimum standard for quality. Alternative game objectives may be included that are accomplished by exceeding a virtual weld performance average over the various levels. Consequently some levels of play may be performed below the performance minimums with others commensurately above. The game objective is met as long as the weighted average for the entire game exceeds a predetermined minimum. [0043] In judging the end user's 11 performance, the simulating device 10 may track the movements of the end users 11 through the input device 13 and compare the data with parameters stored in memory, or coded into the computer program product. The data and/or parameters may be stored in a database, or by any means of data storage chosen with sound judgment. In one embodiment, the simulating device 10 records and stores information about the end user's 11 performance for comparison at a time subsequent to the virtual activity. In other embodiments, comparison with the welding quality parameters is performed in real time with the results being dynamically displayed or catalogued for subsequent review. In addition to the data collected via the input device 13 , other types of data may be captured, which include: time and date data, user name, scenario, as well as game status data. It will be appreciated that any type of data may be tracked and stored as needed for determining and reporting the results of game play. [0044] With reference now to FIG. 7 , as mentioned above, the simulating device 10 may also comprise a system that facilitates training skills used in industrial or commercial settings. In one exemplary embodiment, the simulating device 10 may depict a virtual welding environment 9 a featuring scenarios having one or more themed articles 16 a consistent with that described above. The simulating device 10 may present the end user 11 with a problem to be solved, e.g. a building structure that needs assembled or a race car frame that needs repaired. The problem may be expressly stated whereby the end user 11 is directly tasked with solving the problem given a set of virtual tools 26 . Instructions may be presented describing how the problem should be fixed including which welding techniques or processes should be used. The simulating device 10 may also indicate welder settings or ranges of settings that are needed to weld the virtual article(s) 16 for the given scenario. Additionally, the simulating device 10 may indicate what type of electrode is needed for a particular repair, and/or at what travel speed the weld should be made to make an acceptable virtual weld, which may correspond to a real world weld. However, any type of instruction may be presented to the end user 11 for assembling or repairing the virtual articles 16 . It is once again noted that instructions may be displayed in text or audibly presented in any one of various languages, and/or graphically displayed with graphics as is appropriate for different training settings. [0045] As the end user 11 advances, the level of instruction may be adjusted accordingly. At beginner levels, the level of instruction may focus on fundamentals relating to, for example, welding theory, basic welding practices and/or welder set up. Other training levels may provide tutorials related to various weld joint configurations and/or welding with different types of materials and electrodes. More advanced levels may concentrate on particular welding processes and techniques. Of course, each level may be enhanced by one or more scenarios simulating real world activity as described above. [0046] In one embodiment, the welding training may include problem-based scenarios. The problem-based scenario may be characterized by incorporating an operational deficiency in a virtual article 16 that must be discovered, analyzed, and a solution formulated by the end user 11 . Knowledge learned from a previous lesson or level of training may be relied on for solving the problem. In one example, a race car may be depicted and described as not functioning properly. The virtual environment 9 a may be programmed to present visual, and/or audible, clues that allow the end user 11 to discern the particular problem presented for the given scenario. After analyzing the problem, the end user 11 is directed to devise a solution that, in an exemplary manner, may incorporate: selecting the appropriate welding process, adjusting the welding power supply settings, choosing a particular electrode and then performing a virtual weld. A proper repair therefore requires not only the physical motion of implementing a suitable virtual weld, but also selecting the appropriate welding process and associated parameters. A successful repair or assembly may be indicated, whereby the virtual race car drives away or drives in a race. If an improper or incomplete repair has been made, the race car may perform poorly or not at all with further clues provided to the end user 11 as to what problems remain that need to be fixed. In this manner, welding training encompasses not only the training of muscle memory to successfully perform a particular weld, but also teaches the end user 11 how to properly analyze the virtual article(s) 16 for selecting the appropriate welding process needed to correct its operational deficiency. Welding training may also encompass learning that extends beyond the training of muscle memory by incorporating weld puddle modeling that teaches the end user 11 to make adjustments during the welding process. [0047] As mentioned above, a grade may be derived from the end user's analysis of the problem-based scenario. In one embodiment, the end user 11 may be given information regarding the virtual article's 16 a base material and instructed to select an electrode appropriate for use with that base material. In the real world, selection of an electrode affects the integrity of a weld joint. Similarly selecting the right electrode in the virtual welding environment 9 a affects the score or grade of the end user's 11 performance. Additionally, the end user 11 may be required to calculate the heat input to ensure that the base material properties are not permanently altered by multi-pass welds. In another embodiment, the simulating device 10 may provide the end user 11 with information related to material thickness and/or joint configuration. Accordingly, the end user 11 may be required to determine the appropriate travel speed for the virtual welding power supply settings selected in order to properly make the virtual weld. It is noted here that the information may be expressly stated or indicated by virtual cues from which the end user 11 may infer important factors needed for analyzing the problem. A combination of the aforementioned is also contemplated by the embodiments of the subject invention. It will be recognized that the simulating device 10 therefore functions to educate and evaluate proficiency in learning for science, technology, engineering and/or math as promoted by various educational and governmental institutions. [0048] It may be required that each level of training must be satisfactorily completed before advancing to subsequent levels. In one embodiment, tests may be given related to both welding knowledge and/or virtual welding performance. Data, i.e. test data or performance data, from the current scenario may be tracked, stored and compared against preprogrammed welding parameters in a manner consistent with that described above. In areas where minimum levels of achievement have not been reached, the end user 11 may be given opportunity to review tutorials and/or practice welding a particular weld joint. Once proficiency has been demonstrated, the end user 11 may advance to progressively more difficult levels teaching new skills. Tablet-Based Virtual Welding [0049] Various elements, features, and functions described herein may be embodied in a tablet-based computing device. A tablet-based computing device, or tablet, is generally a mobile, one-piece device having a touch-screen display that a user may interact with via the user's finger or a stylus. Gestures of the user's finger or the stylus serve as a primary means of control and input. However, the tablet may provide additional means of control and input such as, for example, buttons, a virtual or attachable keyboard, or input from one or more sensors. For example, a stylus may have one or more accelerometers and serve as the input device 13 , in accordance with an embodiment. [0050] One embodiment provides a tablet-based computing device or tablet. The terms “tablet-based computing device” and “tablet” are used interchangeably herein. The tablet-based computing device includes a touch-screen display and computer memory storing at least one welding software application providing a virtual welding process. The tablet-based computing device further includes processing means operable to execute coded instructions of the at least one welding software application to generate an interactive virtual welding environment and to display the interactive virtual welding environment on the touch-screen display. The tablet-based computing device also includes an input means configured to interact with the touch-screen display when manipulated by a user to direct at least a spatial orientation of a virtual welding tool in the virtual welding environment while performing a virtual welding activity corresponding to the virtual welding process. The virtual welding process may include one of flux cored arc welding (FCAW), gas metal arc welding (GMAW), gas tungsten arc welding (GTAW), and shielded metal arc welding (SMAW). The input means may include a stylus that does not have a motion sensor or a position sensor. The input means may be configured to interact with the touch-screen display when manipulated by a user to direct at least a spatial position of a virtual welding tool in the virtual welding environment while performing a virtual welding activity corresponding to the virtual welding process. In accordance with an alternative embodiment, instead of having a separate input means, the touch-screen display is configured to be manipulated by a finger of a user to direct at least a spatial orientation of a virtual welding tool in the virtual welding environment while performing a virtual welding activity corresponding to the virtual welding process. Furthermore, the touch-screen display may be configured to be manipulated by a finger of a user to direct at least a spatial position of a virtual welding tool in the virtual welding environment while performing a virtual welding activity corresponding to the virtual welding process. The virtual welding environment may include a virtual asset to be welded using the virtual reality welding process. The virtual asset may be, for example, an automobile, a bridge, a wind turbine tower, or a building. [0051] FIG. 8 illustrates an example embodiment of a tablet-based computing device 800 (or tablet) having a stylus 810 . The tablet 800 also includes a touch-screen display 820 . The stylus 810 may be manipulated by a user to interact with the touch-screen display 820 . For example, the stylus 810 may be manipulated by a user to direct a position and spatial orientation (e.g., angle) of a virtual welding tool, in accordance with an embodiment, as discussed later herein. FIG. 9 illustrates an example embodiment of the tablet-based computing device 800 (or tablet) without the stylus 810 . Instead, a user may use his finger to interact with the touch-screen display 820 . [0052] FIG. 10 illustrates a schematic block diagram of an example embodiment of the tablet 800 . As described previously herein, the tablet 800 includes a touch-screen display 820 and, optionally, a stylus 810 . The tablet 800 also includes a processor 830 and computer memory 840 . The processor 830 may be a programmable microprocessor, for example, although other types of logic processors are possible as well. The computer memory 840 may be, for example, electronic memory, such as a combination of random access memory (RAM) and read-only memory (ROM). Other types of computer memory may be possible as well, in accordance with various other embodiments. [0053] The touch-screen display 820 and the computer memory 840 are operatively connected to the processor 830 . In accordance with an embodiment, the computer memory 840 stores a welding software application (WSA) 850 that provides a virtual welding process (e.g., a flux cored arc welding (FCAW) process, a gas metal arc welding (GMAW) process, a gas tungsten arc welding (GTAW) process, or a shielded metal arc welding (SMAW) process. The WSA 850 provides coded instructions that may be executed on the processor 830 . [0054] In accordance with an embodiment, the WSA 850 provides coded instructions to generate and display an interactive virtual welding environment on the tablet 800 . The virtual welding environment may include a virtual welding tool and a virtual part to be welded, along with various other virtual entities such as, for example, controls, indicators, and assets. The type of virtual welding tool, virtual part to be welded, and various other virtual entities in the virtual welding environment may be dependent upon the virtual welding process. For example, if the WSA 850 corresponds to a SMAW process, then the virtual welding tool will be a stick welding tool. [0055] The virtual controls may allow the user to change a set up of the virtual welding process (e.g., amperage, voltage, wire feed speed, etc.), for example. The virtual indicators may provide displayed indications to the user with respect to, for example, pitch and roll angles of the virtual welding tool. Other indicators are possible as well, as known to one skilled in the art. The virtual assets may correspond to the type of entity that is being welded (e.g., a car, a boat, an air plane, a tractor, a bridge, a wind turbine, or a building). [0056] In accordance with an embodiment, even though a WSA corresponds to a particular welding process, the WSA may allow the selection of different assets in the virtual welding environment. Therefore, for a WSA providing a SMAW welding process, the user may be able to select to weld on a part of a car, a boat, an air plane, a tractor, a bridge, a wind turbine, or a building, for example. Other types of virtual assets are possible as well, in accordance with various other embodiments. [0057] The user may use the stylus 810 or his finger to interact with the touch-screen display 820 to perform a virtual welding activity on the tablet 800 corresponding to the virtual welding process provided by the WSA 850 . For example, the user may direct the position and orientation of a virtual welding tool, displayed on the touch-screen display 820 , with respect to a virtual weld joint to simulate the creation of a weld bead along the weld joint. The user is relying on the interaction of the stylus (or finger) with the touch-screen display to manipulate the displayed virtual weld tool. Therefore, no position sensors or motion sensors are needed to sense the position or orientation of the stylus or the finger. [0058] The computer memory 840 may also store other WSA's corresponding to other virtual welding processes. As a result, a user may select which WSA to execute on the tablet 800 , at any given time. In this manner, a user may have a full complement of virtual welding processes from which to select and practice his virtual welding techniques. [0059] Another embodiment provides a tablet-based, welding computing device or tablet. The tablet includes a display, wireless communication means, and computer memory storing at least one welding summarizing software application. The tablet also includes processing means operable to execute coded instructions of the at least one summarizing software application. The coded instructions are executed to access at least one virtual reality welding system via the wireless communication means to download user virtual welding activity information from the at least one virtual reality welding system to the tablet. The coded instructions are also executed to generate a summary of user virtual welding progress based on the user virtual welding activity information, and display the summary of user virtual welding progress on the display. The user virtual welding activity information and the summary of user virtual welding progress may correspond to a single user of the at least one virtual reality welding system, or to a plurality of users of the at least one virtual reality welding system. [0060] FIG. 11 illustrates an embodiment of a tablet-based computing device 1100 in communication with a virtual reality welding system 1200 via an external communication infrastructure 1300 . The virtual reality welding system 1200 is a simulator used to train welding students how to weld in a virtual environment. As a welding student (user) practices on the virtual reality welding system 1200 , the system 1200 records user virtual welding activity information. The user virtual welding activity information may include, for example, information identifying the types of welding processes the user has performed on the virtual reality welding system 1200 along with information related to a performance of the user for each of the welding processes. The external communication infrastructure 1300 may include one or more of the internet, a cellular telephone network, or a satellite communication network, for example. Other types of external communication infrastructure may be possible as well, in accordance with various other embodiments. [0061] FIG. 12 illustrates a schematic block diagram of an example embodiment of the tablet 1100 of FIG. 11 . The tablet 1100 includes a wireless communication device 1110 . The wireless communication device may include, for example, WiFi communication circuitry and software and/or 3G or 4G communication circuitry and software providing access to the external communication infrastructure 1300 . The tablet 1100 also includes a display 1120 , a processor 1130 , and computer memory 1140 . The display 1120 may be a touch-screen display, in accordance with an embodiment. The processor 1130 may be a programmable microprocessor, for example, although other types of logic processors are possible as well. The computer memory 1140 may be, for example, electronic memory, such as a combination of random access memory (RAM) and read-only memory (ROM). Other types of computer memory may be possible as well, in accordance with various other embodiments. [0062] The computer memory 1140 stores a summarizing software application (SSA) 1145 that provides an information summarizing capability. In accordance with an embodiment, the SSA 1145 provides coded instructions that may be executed on the processor 1130 to access the virtual reality welding system 1200 over the external communication infrastructure 1300 via the wireless communication device 1110 to retrieve user virtual welding activity information stored on the virtual reality welding system 1200 . The user virtual welding activity information may be stored on the virtual reality welding system 1200 in the form of one or more electronic files, for example. [0063] Furthermore, the SSA 1145 provides coded instructions that may be executed on the processor 1130 to generate a summary of user virtual welding progress based on the user virtual welding activity information, and display the summary on the display 1120 . The summary may be displayed on the tablet 1100 in a report format, for example. Other formats are possible as well, in accordance with various other embodiments. [0064] Again, the user virtual welding activity information may include, for example, information identifying the types of welding processes the user has performed on the virtual reality welding system 1200 along with information related to a performance of the user for each of the welding processes. The summary of user virtual welding progress generated by the tablet 1100 may include, for example, average performance information, or consolidated performance information for a user of the virtual reality welding system 1200 . For example, an average pitch angle of how the user held a mock welding tool of the virtual reality welding system 1200 during a particular virtual welding process may be generated by the SSA 1145 and displayed as part of the summary. Furthermore, a consolidated presentation of pitch angle vs. roll angle of how a user held the mock welding tool during a particular virtual welding process may be generated by the SSA 1145 and displayed as part of the summary. [0065] The summary of user virtual welding progress may also include graphical information showing how a performance parameter associated with a user has changed (e.g., improved) over time. For example, a graph of the average travel speed of a mock welding tool over a plurality of successive welding activities performed by a user for a particular welding process may be generated by the SSA 1145 and displayed as part of the summary. The graph may indicate how the average travel speed started out varying between too fast and too slow and then eventually settled to a desired travel speed during the course of, for example, twenty (20) successive welding activities for a particular welding process, thus providing an indication of how long it took for the user to settle into applying the correct travel speed to the mock welding tool of the virtual reality welding system 1200 . [0066] User virtual welding activity information may be accessed for a single user, or for a plurality of users, from one or more virtual reality welding systems, in accordance with an embodiment. For example, a welding instructor, using the tablet 1100 , may access user virtual welding activity information for all of his welding students across a plurality of virtual reality welding systems 1200 . The SSA 1145 on the tablet 1100 may create a summary for each welding student and may also create a consolidated summary which shows progress for all of the welding students, for example, in a comparative manner (e.g., a ranking of the welding students). [0067] A further embodiment provides a method of virtual welding. The method includes generating a dynamic virtual welding environment within a computer-based platform and displaying stereoscopic three-dimensional (3D) images of at least a portion of the dynamic virtual welding environment on a display screen of the computer-based platform. The computer-based platform may be one of a desktop personal computer, a tablet-based computer device, a laptop computer, a notebook computer, or a workstation, for example. The method further includes viewing the stereoscopic three-dimensional images using 3D glasses, resulting in at least a portion of the dynamic virtual welding environment appearing to project out of the display screen into 3D space. The method also includes virtually welding a virtual weldment of the projected portion of the dynamic virtual welding environment in 3D space using a mock welding tool while viewing the stereoscopic three-dimensional images using 3D glasses. In accordance with an embodiment, a position and an orientation of the mock welding tool is tracked in 3D space in real time using one or more of inertial tracking techniques or magnetic tracking techniques, for example. The method also includes calibrating the mock welding tool to the projected virtual welding environment in 3D space by, for example, correlating a position of at least one point on the mock welding tool in 3D space to a position of at least one point on the projected virtual weldment in 3D space. [0068] FIG. 13 illustrates a portion of a dynamic virtual welding environment 1310 , generated by a computer-based platform 1300 , that is stereoscopically projected out of a display screen 1320 of the computer-based platform 1300 into 3D space as viewed by a user wearing 3D glasses, in accordance with an embodiment. The computer-based platform 1300 is configured to simulate welding activity in the virtual welding environment 1310 in real time, as affected by a user. The user may employ a mock welding tool 1330 (e.g., a plastic tool simulating a real-world welding tool) to interact with the projected portion of the dynamic virtual welding environment 1310 to virtually weld a virtual weldment (e.g., a virtual pipe 1311 ) to create a virtual welded joint 1312 appearing to the user in 3D space. The resultant virtual welded joint 1312 includes a virtual weld bead as created during the virtual welding activity. [0069] The virtual welding environment is dynamic in the sense that it may be acted upon and modified in response to the user performing a virtual welding activity in real time. As the user moves the mock welding tool along the joint to create the virtual welded joint 1312 appearing to the user in 3D space, the computer-based platform 1300 updates the dynamic virtual welding environment 1310 in real time and continues to display stereoscopic 3D images of the environment on the display screen 1320 such that the user may observe his actions and the resultant creation of the virtual welded joint 1312 in real time in 3D space using 3D glasses. [0070] In accordance with an embodiment, the appearance of the resultant virtual welded joint 1312 is realistically simulated as being dependent on the user's technique of manipulating the mock welding tool 1330 , as would be the case in the real-world during a real welding activity. For example, if the user moves the tip of the mock welding tool 1330 away from the joint 1312 , the displayed deposited weld bead will appear as being deposited away from the joint 1312 . Furthermore, if the user moves the tip of the mock welding tool 1330 too quickly along the joint 1312 , the resultant displayed deposited weld bead may have a distorted stacked-dime appearance. [0071] For the mock welding tool 1330 to effectively interact with the virtual welding environment, as viewed as projected into 3D space by the user, the mock welding tool is tracked in 3D space and is correlated to the projected virtual welding environment in 3D space by the computer-based platform 1300 . The mock welding tool may be tracked by the computer-based platform 1300 in 3D space by employing, for example, inertial tracking techniques or magnetic tracking techniques. An inertial tracking technique may employ accelerometers on the mock welding tool which report position and motion information back to the computer-based platform 1300 via wired means (e.g., a universal serial bus connection) or wireless means (e.g., a Bluetooth™ connection). A magnetic tracking technique may employ coil sensors on the mock welding tool and a magnetic source providing a magnetic field to be detected by the coil sensors. In a similar manner, the coil sensors may report position and motion information back to the computer-based platform 1300 via wired or wireless means. [0072] To correlate or calibrate the position of the mock welding tool 1330 to the “projected position” of the virtual welding environment 1310 in 3D space, the computer-based platform 1300 may direct the user to successively place the tip of the mock welding tool 1330 at two or more points on an element of the projected environment. For example, the computer-based platform 1300 may direct the user to place the tip of the mock welding tool at pre-determined points “A”, “B”, and “C” on the virtual weldment 1311 in succession (see FIG. 13 ) as the user views the projection of the virtual weldment 1311 in 3D space. Upon recording the positions of the tip of the mock welding tool at points “A”, “B”, and “C” and performing a correlation calculation, the computer-based platform 1300 “knows” where the tip of the mock welding tool is with respect to the user's perception of the projected virtual welding environment in 3D space as the user performs the virtual welding activity. [0073] FIG. 14 is a flowchart of an example embodiment of a method 1400 of virtual welding. In step 1410 , generate a dynamic virtual welding environment within a computer-based platform. In step 1420 , display stereoscopic 3D images of at least a portion of the dynamic virtual welding environment on a display screen of the computer-based platform. In step 1430 , view the stereoscopic 3D images using 3D glasses, resulting in at least a portion of the dynamic virtual welding environment appearing to project out of the display screen into 3D space. In step 1440 , virtually weld a virtual weldment of the projected portion of the dynamic virtual welding environment in 3D space using a mock welding tool while viewing the stereoscopic 3D images using 3D glasses. Again, the computer-based platform may be one of a desktop personal computer, a tablet-based computer device, a laptop computer, a notebook computer, or a workstation, for example. In accordance with an embodiment, the computer-based platform tracks the position and motion of the mock welding tool during the virtual welding activity, where the mock welding tool and the projected virtual welding environment have been correlated in 3D space. [0074] While the claimed subject matter of the present application has been described with reference to certain embodiments, it will be understood by those skilled in the art that various changes may be made and equivalents may be substituted without departing from the scope of the claimed subject matter. In addition, many modifications may be made to adapt a particular situation or material to the teachings of the claimed subject matter without departing from its scope. Therefore, it is intended that the claimed subject matter not be limited to the particular embodiments disclosed, but that the claimed subject matter will include all embodiments falling within the scope of the appended claims.

Embodiments of the present invention pertain to a computer program product and processor based computing system that provides processing means for executing coded instructions and input means for interacting with said processing means to create a virtual welding environment. The system establishes an objective to change a functional or operational state of a virtual article, and directs the end user to perform at least one virtual welding operation for changing its functional state. The system trains new users and inexperienced welders on the fundamental aspects of welding and other technical aspects.

BACKGROUND OF THE INVENTION [0001] This application relates to a latching mechanism. More specifically, but not by way of limitation, this invention relates to a latching mechanism for lifting and lowering containers. [0002] In industrial applications, operators find it necessary to lift containers. Many times, the containers are located in remote areas. For instance, in the offshore energy industry, boats ferry equipment to platforms located many miles from shore. As those of ordinary skill in the art will recognize, the seas can get quite rough. Generally, the containers will be positioned on the aft deck of the boat. In order to offload the equipment, a platform crane is used to lift the containers from the boat deck onto the platform. [0003] Prior art devices requires the use of a ball hook, which is attached to the cable. A deck hand is required to attached the ball hook onto the containers. Sometimes, the ball hook is attached to a set of slings. As those of ordinary skill in the art will recognize, attaching the ball hook onto the slings is a dangerous endeavor since the boat may be rocking due to wind and/or waves. For instance, in 10 foot seas, the boat may be traveling vertically as much as 20 feet since the boat will ride the crest of the wave to the trough of the wave. However, the ball hook is not moving in unison with the boat since the crane is mounted on the platform. Additionally, in some applications, the crane may be mounted on floating platform. Therefore, the potential for injury to the deck hand is high due to the logistics of rigging up and rigging down the ball hook to the container. [0004] Therefore, there is a need to have an apparatus and mechanism that allows quick and safe latching and unlatching onto containers. There is also a need that will allow for the lifting containers. There is also a need for a device and method that allows for containers on vessels to be safely loaded and/or offloaded. These needs, and many others, will be met by a reading of the following detailed description of the preferred embodiments. SUMMARY OF THE INVENTION [0005] A latching apparatus for lifting and lowering loads is disclosed. In the preferred embodiment, the apparatus comprises a prong having a protuberance at a first end and a receiving receptacle configured to receive the first end of the prong. In one preferred embodiment, the receiving receptacle comprises a body having a central passage, a plurality of jaws arranged about the central passage, and biasing means for biasing the jaws to extend into the central passage and engaging the protuberance. The biasing means may comprise a spring having a first end engaging the jaws and a second end engaging an internal portion of the body so that the jaws are biased toward the central passage. [0006] The latching apparatus may further comprise an activation disc for disengaging the jaws from the protuberance. In one preferred embodiment, the jaws have a pin extending therefrom; and wherein the activation disc contains a plurality of cam surfaces that engage the pin of the jaws. Additionally, the body may comprise a first section that is attached to a second section, and wherein the second section contains the jaws and the activation disc, and wherein the first section is configured to abut and align the plurality of jaws. [0007] The latching apparatus may further comprise a handle mounted through the body so that the handle has an exterior portion and an internal portion, and wherein the handle has a projection that engages the activation disc so that movement of handle creates movement of the activation disc. The latching apparatus may further comprise an alignment skirt attached to the body, and wherein the alignment skirt has a conical shape, and wherein the alignment skirt is configured to direct the prong into the central passage. [0008] In one embodiment, the spring is disposed about an indicator pin, and the indicator pin extends through an aperture in the body when the jaw inserts are disengaged from the protuberance. Also, the latching mechanism may further comprise a jaw insertion device for engaging the jaw inserts from the aperture within the body. [0009] A method of latching onto a container is also disclosed. The method includes providing a prong having a protuberance at a first end, and wherein the prong is attached to the container, and wherein the protuberance has a neck portion formed thereon. The method further comprises providing a receiving receptacle configured to receive the first end of the prong, the receiving receptacle comprising: a body having a central passage; a plurality of jaws arranged about the central passage; and, a spring having a first end engaging the jaws and a second end engaging an internal portion of the body so that the jaws are biased into the central passage. The method further includes lowering the receiving receptacle so that the protuberance enters the central passage, engaging the protuberance with the jaws, and compressing the spring. The method further includes opening the jaws while continuing to lower the receiving receptacle, allowing the jaws to close about the neck portion, and latching the protuberance with the jaws. [0010] In one preferred embodiment, the receiving receptacle has attached a cable, and the method further comprises pulling on the cable to lift the container. Also, in one preferred embodiment, the receiving receptacle further comprises an activation disc for disengaging the jaws from the protuberance, and wherein the method further comprises turning a handle which rotates a projection on the handle and engaging the projection with the activation disc so that rotation of the projection causes rotation of the activation disc. Next, the activation disc is rotated and the pin on the jaws is engaged with a cam surface on the activation disc via the rotation of the activation disc. Next, the spring is compressed, and the jaws are expanded and released from the protuberance, thereby unlatching the receiving receptacle from the prong. [0011] Additionally, in one preferred embodiment, the prong has a cylindrical body and the first end of the jaws is curved, and wherein the step of latching the protuberance with the jaws includes engaging the curved end of the jaws about the cylindrical body of the prong. Also, a funnel may extend from the body, and wherein the step of lowering the receiving receptacle includes guiding the jaws onto the prong with the funnel. [0012] An advantage of the present invention is the ability to connect and disconnect quickly and safely. Another advantage is that the present apparatus can offload equipment from floating vessels. Still yet another advantage is that the receptacle can be latched onto the prong without the need to have a person directly operating the latch. In other words, the apparatus can be lowered from a platform on a cable, and the prong can be attached to a container on a boat deck, and due to the design, the receptacle can be latched onto the prong by having the crane operator lower the skirt onto the prong without having an individual on the deck of the boat. Still yet another advantage is the ability to remotely control the handle means in order to turn the handle [0013] A feature of the present invention is a fishing neck type of prong that is configured to engage with the jaw mechanism. A feature of the jaw mechanism is the multiple inserts that are used to engage with the prong. Another feature is the jaw mechanism is a spring loaded latch mechanism which allows entry of the prong, and thereafter, captures the prong. Yet another feature is the disc activation device that is used to unlatch the apparatus from the prong. Still yet another feature is the skirt that serves to funnel the prong into the jaw mechanism. BRIEF DESCRIPTION OF THE DRAWINGS [0014] FIG. 1 is a perspective view of the latching mechanism as the receiving receptacle is being lowered. [0015] FIG. 2 is an exploded view of the most preferred embodiment latching mechanism that includes the receiving receptacle and the prong. [0016] FIG. 3 is a partial cross-sectional view of the jaws and activation disc of the most preferred embodiment. [0017] FIG. 4 is a partial cross-sectional view of the activation disc seen in FIG. 3 . [0018] FIG. 5 is an expanded view of the cam surface of the activation disc seen in FIG. 4 . [0019] FIG. 6 is a cross-sectional view of the latching mechanism taken along line 7 - 7 of FIG. 7 . [0020] FIG. 7 is a partial cut-away view of the handle means of the latching mechanism corresponding to the jaws being in the contracted position. [0021] FIG. 8 is a sequential partial cross-sectional view of the jaws and activation disc with the jaws being in the expanded, open position. [0022] FIG. 9 is a partial cross-sectional view of the actiavation disc seen in FIG. 8 . [0023] FIG. 10 is an enlarged view of a cam surface seen in FIG. 9 . [0024] FIG. 11 is a cross-sectional view of the latching mechanism seen in FIG. 6 with the jaws in the expanded, open position. [0025] FIG. 12 is a partial cross-sectional view of the handle means of the latching mechanism corresponding to the jaws being in the expanded position. [0026] FIG. 13 is a partial cross-sectional view of the jaws and activation disc, with the jaws about the prong, and depicting the jaw insertion device. [0027] FIG. 14 is a perspective view of the latching mechanism and the receiving receptacle latched together. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS [0028] FIG. 1 illustrates a perspective view of the most preferred embodiment of the latching mechanism 2 that includes the receiving receptacle 4 and the prong 6 . FIG. 1 further illustrates the support members 8 , 10 , 12 that will be connected to a container 15 . The receiving receptacle 4 is being lowered via cable 14 that is connected to a shackle. [0029] Referring now to FIG. 2 , an exploded view of the most preferred embodiment of the latching mechanism 2 will now be described. FIG. 2 depicts the body 16 , wherein the body includes a flange end 18 that extends to a cylindrical body 20 that in turn extends to the conically shaped outer surface 21 which in turn terminates at the radial end 22 . Extending radially inward is the conically shaped inner surface 24 , and wherein the conically shaped inner surface 24 may be referred to as the funnel 24 . The funnel 24 extends to the central bore passage 26 . [0030] As seen in FIG. 2 , the flanged end 18 extends radially inward to a groove section 28 . The groove section 28 is configured to have a cylindrical inner wall 30 , and wherein three radial grooves are formed therein (only radial grooves 32 , 34 are shown in FIG. 2 ). An activation disc 36 is shown, and wherein the activation disc 36 contains the cam surfaces 38 , 40 , 42 . The activation disc 36 is configured to be positioned within the groove section 28 . FIG. 2 further depicts the jaw inserts 44 , 46 , 48 , wherein the jaws 44 , 46 , 48 are configured to be positioned within the radial grooves (i.e. radial grooves 32 , 34 ). The jaws 44 , 46 , 48 have a first end that has a curved portion such as seen at 50 . The jaws 44 , 46 , 48 also contain a pin on the underside, such as pins 52 , 54 , and wherein the pins will engage the cam surfaces 36 , 38 , 40 of the activation disc 36 as will be more fully described. [0031] FIG. 2 also depicts the cover plate 56 to cooperate and engage with the body 16 , and more specifically, the flange end 18 . As seen in FIG. 2 , nuts and bolts (such as nut 58 and bolt 60 ) are used to attach the cover plate 56 to the body 16 . It should be noted that the cover plate 56 is configured to abut and align the top side of the plurality of jaws and activation disc. FIG. 2 also shows the eyelet plate 62 that can be used to attach to a cable via a shackle, as well understood by those of ordinary skill in the art. FIG. 2 also depicts the handle 64 that is disposed through the body 16 . The handle 64 has a radially extending arm 66 , and wherein a projection 68 extends perpendicularly from the arm 66 . The projection 68 will fit through the opening 70 that is disposed through the activation disc 36 . Hence, the lateral movement of the projection 68 (by virtue of the handle 64 being turned) will in turn cause the activation disc 36 to rotate, as will be more fully explained later in the application. [0032] The prong 6 is also shown in FIG. 2 . More specifically, the prong 6 is generally a cylindrical member that contains a protuberance 72 that extends to a neck portion 74 , and wherein the neck portion 74 is a reduced diameter groove about the stem 76 . The protuberance 72 contains the chamfered surface 78 . The jaws 44 , 46 , 48 will cooperate and engage with the neck portion 74 in order to latch the receiving receptacle with the prong 6 . FIG. 2 depicts the spring 82 that will be disposed about indicator pin 84 ; spring 88 that will be disposed about indicator pin 90 ; and, spring 94 that will be disposed about an indicator pin (not shown in this view), as will be more fully set out below. [0033] This embodiment also depicts the jaw insertion device 104 . In the event that the jaws 44 , 46 , 48 become stuck, it is within the teachings of this invention that the indicator pins 84 , 90 , 96 , operatively associated with each jaw, be removed. In order to remove the indicator pins 84 , 90 , 96 , the operator could simply unscrew the indicator pins from the jaws. After the indicator pins 84 , 90 , 96 are removed, the jaw insertion device 104 is threaded into the respective jaw. In the most preferred embodiment, the jaw insertion device 104 has an external threaded end that will engage an internal thread means within the jaw. Hence, after engaging the jaw insertion device 104 , the operator can exert a pull force on the device 104 which would in turn move the jaw 44 thereby unsticking and/or dislodging the jaw. [0034] Referring now to FIG. 3 , a partial cross-sectional view of the jaws 44 , 46 , 48 and activation disc 36 of the most preferred embodiment will now be described. In the view of FIG. 3 , the cover plate 56 has been removed. FIG. 3 depicts the spring 82 that is disposed about the indicator pin 84 , and wherein indicator pin 84 has a first end within the aperture 86 and a second end threadedly engaged within the jaw 44 . As shown in FIG. 3 , the spring 82 biases the jaw 44 toward the axial center i.e. against the neck portion 74 of the prong 6 . [0035] FIG. 3 also depicts the spring 88 that is disposed about the indicator pin 90 , and wherein pin 90 has a first end disposed within the aperture 92 and a second end threadedly engaged within the jaw 46 , wherein the spring 88 biases the jaw 46 toward the axial center i.e. against the neck portion 74 of the prong 6 . The spring 94 that is disposed about the indicator pin 96 , and wherein indicator pin 96 has a first end within the aperture 98 and a second end threadedly engaged within the jaw 48 , and the spring 94 biases the jaw 48 toward the axial center i.e. against the neck portion 74 of the prong 6 . [0036] In order to disengage the jaws 44 , 46 , 48 from the neck portion of prong 6 , the activation disc 36 will have to be rotated via the projection 68 . Referring now to FIG. 4 , a partial cross-sectional view of the activation disc 36 seen in FIG. 2 will now be described. More specifically, FIG. 4 shows the top view of the activation disc 36 without the jaws 44 , 46 , 48 . FIG. 4 shows the position of the pin 52 for jaw 44 , pin 54 for jaw 48 and pin 100 for jaw 46 . Hence, pin 52 is abutting cam surface 42 , pin 100 is abutting cam surface 38 and pin 54 is abutting cam surface 40 . As seen, the cam surface 38 , 40 , 42 are curved notches within disc 36 . FIG. 5 is an expanded view of the cam surface 42 of the activation disc 36 seen in FIG. 4 . As will be more fully set out in the application, as the activation disc 36 is rotated, the sloping cam surface 42 will cause the pin 52 (and in turn the jaw 44 ) to be drawn towards the inner wall 101 a thereby opening the jaws. A stop pin 101 b is shown which stops movement of the activation disc 36 beyond a predetermined point. [0037] Referring now to FIG. 6 , a cross-sectional view of the latching mechanism 2 taken along line 6 - 6 of FIG. 7 will now be described. This view depicts the jaws engaging the neck portion, and in particular the jaw 46 engaging neck portion 74 , and wherein the spring 88 has biased the jaw 46 into engagement. Hence, the prong 6 is engaged. The handle 64 and the radially connected arm 66 is shown, and wherein the projection 68 is shown disposed through the activation disc 36 . Note that the indicator pin 90 is not visible in the position shown in FIG. 6 which indicates the jaws engaged position. [0038] A partial cut-away view of the handle 64 of the latching mechanism 2 corresponding to the jaws being in the contracted position is shown in FIG. 7 . The partial cut-away within the cylindrical body 20 depicts the projection 68 disposed through the activation disc 36 . FIG. 7 further shows the line 3 - 3 for the partial cross-sectional view of FIG. 3 and line 4 - 4 for the partial cross-sectional view of FIG. 4 . [0039] Referring now to FIG. 8 , a sequential partial cross-sectional view of the jaws 44 , 46 , 48 and activation disc 36 with the jaws being in the expanded open position will now be described. More specifically, the handle 64 has been rotated which in turn causes the projection 68 to rotate the activation disc 36 . The pins 52 , 54 , 100 (not seen in this view) follow the cam surfaces 38 , 40 , 42 which in turn cause the jaws 44 , 46 , 48 to expand (open). As seen in FIG. 8 , the springs 82 , 88 , and 94 are compressed do to this movement. [0040] FIG. 8 additionally depicts the feature of the indicator pins 84 , 90 , 96 disposed through the apertures. As seen in FIG. 8 , the indicator pins 84 , 90 , 96 serve as indicator means for indicating the position of the jaws. In other words, the visible ends of indicator pins 84 , 90 , 96 indicate that the jaws are in the open position. In this way, an operator can visually determine the position of the jaws and whether the prong is latched or unlatched. Hence, when the operator sees the indicator pin, the operator can determine that the jaw is not latched with the prong. [0041] FIG. 9 is a partial cross-sectional view of the activation disc 36 seen in FIG. 8 . FIG. 9 depicts the pins 52 , 54 , 100 in relation to the cam surfaces 42 , 40 , 38 . FIG. 10 is an enlarged view of a cam surface 42 seen in FIG. 9 engaging the pin 52 . As seen in FIG. 10 , the movement of the cam surface 42 has caused the pin 52 (and in turn jaw 44 ) to move towards inner wall 101 a. [0042] Referring now to FIG. 11 , a cross-sectional view of the latching mechanism 2 seen in FIG. 6 with the jaws in the expanded (open) position. Hence, the projection 68 has caused the rotation of the activation disc 36 , which in turn caused the jaw to expand thereby contracting the spring 88 . As seen in FIG. 11 , jaw 46 is no longer engaging the neck portion 74 . Also, the indicator pin 90 is extending outward from the body 20 indicating an open jaw. [0043] FIG. 12 is a partial cut-away view of the handle 64 of the latching mechanism 2 corresponding to the jaws being in the expanded position. The projection 68 is shown disposed through the activation disc 36 . It should also be noted that the handle may be manually, hydraulically and/or pneumatically controlled. Also, it is within the teachings of this invention that actuation of the handle 64 can be done remotely i.e. remote control of the handle in order to turn the handle. FIG. 12 also shows the line 8 - 8 for the partial cross-sectional view of FIG. 8 , the line 9 - 9 for the partial cross-sectional view of FIG. 9 , and the line 11 - 11 for the partial cross-sectional view of FIG. 11 . [0044] FIG. 13 depicts a partial cross-sectional view of the jaws 44 , 46 , 48 and activation disc 36 , with the jaws about the prong, and depicting the jaw insertion device. FIG. 14 is a perspective view of the latching mechanism and the receiving receptacle latched together. [0045] Although the present invention has been described in terms of specific embodiments, it is anticipated that alterations and modifications thereof will no doubt become apparent to those skilled in the art. It is therefore intended that the following claims be interpreted as covering all such alterations and modifications as fall within the true spirit and scope of the invention.

A latching apparatus for lifting and lowering loads. The apparatus comprises a prong and a receiving receptacle configured to receive the prong. The receiving receptacle comprises a body having a central passage, a plurality of jaws arranged about the central passage, and a spring for biasing the jaws to extend into the central passage. The latching apparatus may further comprise an activation disc for disengaging the jaws from the prong. The jaws have a pin extending therefrom; and wherein the activation disc contains a plurality of cam surfaces that engage the pin of the jaws. A method of latching onto a container is also disclosed.

RELATED APPLICATIONS This application claims the benefit of U.S. Provisional Application No. 60/018,986, filed Jun. 4, 1996, and of U.S. Provisional Application No. 60/042,602, filed Apr. 2, 1997. This application is also related to U.S. Pat. No. 5,612,597, issued Mar. 18, 1997, entitled "OSCILLATING DRIVER CIRCUIT WITH POWER FACTOR CORRECTION, ELECTRONIC LAMP BALLAST EMPLOYING SAME AND DRIVER METHOD" (IR-1166), to U.S. Pat. No. 5,545,955, issued Aug. 13, 1996, entitled "MOS GATE DRIVER FOR BALLAST CIRCUITS" (IR-1074), to U.S. Pat. No. 5,559,394, issued Sep. 24, 1996, entitled "MOS GATE DRIVER FOR BALLAST CIRCUITS" (IR-1252), all of which are in the name of the present inventor, and to U.S. Pat. No. 5,550,436, issued Aug. 27, 1996, entitled "MOS GATE DRIVER INTEGRATED CIRCUIT FOR BALLAST CIRCUITS" (IR-1055), in the name of Talbott M. Houk. All of the above are assigned to the assignee of the present application. BACKGROUND OF THE INVENTION The present invention is directed to a lamp ballast circuit which includes a gate driver circuit for driving MOS gated devices and, more particularly, to a lamp ballast circuit which is protected against the removal of or failure of the lamp. Electronic ballast circuits for driving fluorescent lamps or other gas discharge illumination devices are coming into widespread use because of the availability of power MOSFET switching devices to replace the previously used bipolar transistor devices. Typically, the electronic ballast circuit uses two power MOSFET switches in a totem pole (half-bridge) arrangement which includes one or more L-C series resonance circuits and in which the lamp or lamps are connected across one of the reactances of the L-C circuit. The power MOSFET switches are driven to conduct alternately by inputs from the secondary windings of a current transformer whose primary winding conducts the current of the lamp circuits. The primary winding current alternates at the resonance frequency of the resonant circuit. Recently, integrated circuit MOS gate driver devices have been introduced in place of the current transformers. These integrated circuit devices drive the power MOSFETs or IGBTs of an inverter circuit in the ballast circuit from logic level ground referenced inputs and provide a self-oscillating function which is particularly suited for use in electronic lamp ballast circuits. The integrated circuit devices significantly save cost, weight and space when compared to driver circuits employing current transformers. The MOS gate driver circuits have the drawback, however, that when the lamps are removed or are broken and create an open circuit, the power MOSFETs continue to switch in a hard switching mode with high power dissipation and fail due to overtemperature. It is therefore greatly desirable that the ballast circuit be able to shut down when there is an open circuit condition. SUMMARY OF THE INVENTION The present invention provides for an electronic lamp ballast circuit which, when a lamp is removed or fails, is shut down. In accordance with an aspect of the present invention, the operating voltage is supplied to an MOS gate driver circuit by one or more secondary windings of a transformer/inductor having corresponding primary windings that are each in series with one of the lamps. When the lamps are removed or fail, resulting in an open load condition, current is interrupted to the primary side of the transformer/inductor, and, as a result, no current is delivered by the secondary windings which, in turn, reduces the operating voltage supplied to the MOS gate driver circuit to less than the undervoltage lockout voltage. According to another aspect of the present invention, the lamps are coupled to the gate of a switch, such as a power MOSFET, whose source and drain are connected between the operating voltage input of the MOS gate driver circuit and the ground terminal. When the lamps are removed or fail, the switch is activated and clamps the operating voltage input to the ground terminal, hereby shutting off the gate driver circuit. According to a further aspect of the invention, a soft-starting circuit gradually increases the voltage across the lamps prior to their ignition. A PTC thermistor or an L-C circuit gradually heats the lamps until the voltage across the lamps is sufficient for the lamps to strike. Other objects, features and advantages of the present invention will be apparent from the detailed description below. BRIEF DESCRIPTION OF THE DRAWINGS The invention will now be described in greater detail in the following detailed description with reference to the drawings in which: FIG. 1 is a schematic diagram showing an electronic ballast circuit with lamp removal protection according to an embodiment of the present invention; FIG. 2 is a schematic diagram showing an electronic ballast circuit with lamp removal protection according to another embodiment of the present invention; FIG. 3 is a schematic diagram showing an electronic ballast circuit with lamp removal protection according to still another embodiment of the present invention; FIG. 4 is a schematic diagram showing a further embodiment of an electronic lamp ballast circuit with lamp removal protection according to the invention; and FIG. 5 shows a schematic diagram of a still further embodiment of an electronic ballast circuit with lamp removal protection according to the invention. DETAILED DESCRIPTION OF THE INVENTION FIG. 1 shows the novel ballast circuit according to an embodiment of the present invention. More specifically, the circuit of FIG. 1 includes a gas discharge lamp 10 associated with an L-C circuit formed of capacitors 44 and 72 and primary winding inductor 62b of transformer 62 and a second lamp 12 associated with an L-C circuit formed of capacitors 44 and 70 and primary winding inductor 60b of transformer 60. The circuit is typically connected to the lamps using a standard "double 40" fixture. Two power MOSFETs 20 and 21 are connected to dc source terminals V+ in a totem-pole or half-bridge configuration. Power MOSFETs 20 and 21 may be any power device which has a MOS gate, for example, an IGBT or a MOS gated thyristor. The MOS gate driver chip 30 of FIG. 1 provides drive signals to the MOSFETs 20 and 21. More specifically, the chip 30 may be housed in an 8-pin DIP or surface mount package, and has the following pinouts: H O --an output pin to the gate of the high side MOSFET 20. L O --an output pin to the gate of the low side MOSFET 21. V S --a pin to the center tap of the totem-pole or half-bridge connected MOSFETs 20 and 21. G--a pin connected to the negative or ground terminal 23 of the d-c source. C T --a single input control pin which is connected to the node between timing capacitor 38 and timing resistor 36. The other side of capacitor 38 is connected to the negative or ground terminal. The signal at pin C T controls both outputs H O and L O . R T --a pin which is connected to the other terminal of timing resistor 36. V CC --a pin which receives a chip operating voltage from the secondary windings 60a and 62a. V B --a pin connected to the node of diode 32 and capacitor 34, which acts as a "bootstrap" circuit to provide power for the operation of the high side switch. In operation, and before tube 10 strikes, for example, the resonant circuit for this lamp consists of inductor 62b and both of capacitors 44 and 72. The capacitance of capacitor 72 is lower than that of capacitor 44 so that it operates at a higher a-c voltage than that of capacitor 44. This voltage on capacitor 72 strikes the lamp 10. After lamp 10 strikes, capacitor 72 is effectively shunted by the lamp voltage drop, and the frequency of the resonant lamp circuit now depends on inductor 62b and capacitor 44. The shunting of capacitor 72 causes a shift to a lower resonant frequency during normal operation. The voltage of capacitor 72 is rectified by diode 68 and supplies current via resistor 64 to raise the supply voltage of chip 30 above the undervoltage lockout voltage to initiate oscillation. As noted above, the oscillation frequency of the circuit is synchronized by the resistor 36 and the capacitor 38. Chip 30 also has interior circuitry to provide a dead time between outputs of the alternating high side and low side outputs for driving switches 20 and 21. This dead time could vary depending upon the particular application of the I.C. to: 1.) prevent cross-conduction currents from flowing in the power MOSFETs 20 and 21, and 2.) allow an external "snubber" circuit, such as resistor 41 and capacitor 54, to control the half-bridge output voltage slew rate in order to reduce radiated EMI noise. Also, the chip 30 may be supplied at terminal V CC with a rectified a-c voltage. The voltage supplied at terminal V+ can range from 225 volts d-c to about 450 volts d-c, depending upon the supplied AC input voltage or the supplied PFC output voltage. The inductance values of inductors 60b and 62b will depend on the value of voltage V+ and are selected so that the oscillation frequency of the circuit ranges from approximately 30 kHZ to about 45 kHZ. The ballast circuit of this embodiment of the present invention is capable of operating with either one lamp or two lamps inserted in the fixture. Alternatively, the circuit may be configured to operate with one, two or more lamps by providing additional lamp circuits that are each connected to a respective capacitor, transformer and diodes arranged in parallel with and having values similar to that of diodes 66 and 68, capacitors 70 and 72, transformers 60 and 62 and diodes 46 and 48 of lamps 10 and 12, respectively. The novel protection feature of the present invention operates as follows: When both of the two lamps 10 and 12 are removed or fail, the two current paths that include lamps 10 and 12, respectively, and primary windings 60b and 62b of transformers 60 and 62, respectively, are interrupted so that no current flows through the primary windings 60b and 62b. As a result, the current flow through the secondary windings 60a and 62a of transformers 60 and 62 is interrupted, and reduced operating voltage is supplied to inputs V CC and V B of the MOS gate driver chip 30. The gate driver chip 30 therefore does not supply voltage or current to outputs H O and L O so that the power MOSFETs 20 and 21 are turned off and do not supply power to the lamps 10 and 12. It should be noted that because V CC is generally lowered below the undervoltage lockout voltage, the circuit does not "hiccup" when both lamps fail or are removed. When at least one of lamps 10 or 20 is subsequently inserted in the fixture, the current path that includes its primary winding 60b or 62b of transformer 60 or 62 is closed so that an operating voltage is supplied by the corresponding secondary winding 60a or 62a to inputs V CC and V B . The gate driver chip 30 is then activated which, in turn, activates the lamp. FIG. 2 illustrates another embodiment according to the present invention in which PTC thermistors 90 and 92 are used to preheat the lamps 10 and 12. The elements shown in FIG. 2 that have the same reference numerals as are shown in FIG. 1 are identical to the corresponding elements shown in FIG. 1. The lamp circuit in FIG. 2 likewise employs two 40 watt fluorescent lamps 10 and 12 in a common reflector which have respective series inductors 62a and 60a and series capacitor 44. Each of tubes 10 and 12 have parallel capacitors 70 and 72, respectively, and parallel positive temperature coefficient thermistors 90 and 92, respectively. Soft-starting with tube filament pre-heating is accomplished by the P.T.C. thermistors 90 and 92 across each lamp. In this way, the voltage across the lamp gradually increases as the P.T.C. thermistor self-heats until finally the striking voltage with hot filaments is reached and the lamp strikes. The circuit shown in FIG. 2 also includes the novel lamp protection feature shown in FIG. 1. FIG. 3 illustrates yet another embodiment of the present invention which includes an alternative implementation of the novel protection feature of the present invention. FIG. 3 shows a ballast circuit in which the input a-c circuit includes a 230 volt a-c source having two a-c terminals 101a and 101b and which supplies power to an EMI filter circuit 100. The filter circuit supplies power to a single phase full wave rectifier 102 having a positive output connected to resistor 105 and a negative terminal connected to capacitor 104 and which provides a 320 volt d-c output from the 230 volt a-c input. The MOS gate driver chip 30 of FIG. 3 operates directly off the d-c bus through the dropping resistor 105 and oscillates in compliance with the following relationship: ##EQU1## where R 136 is the resistance of resistor 136 and C 138 is the capacitance of capacitor 138. Power for the high side switch gate drive comes from bootstrap capacitor 144 which is generally much larger than the MOSFET input capacitance and is charged to approximately 14 volts whenever pin V S is pulled low during the low side power switch conduction. The bootstrap diode 112 blocks the d-c bus voltage when the high side switch conducts. Diode 112 is a fast recovery diode to ensure that the bootstrap capacitor 144 is not partially discharged as the diode 112 recovers and blocks the high voltage bus. The high frequency output from the half-bridge MOSFETs 20 and 21 would ordinarily be a square wave with very fast transition times. In order to avoid excessive radiated noise from the fast wave fronts, a snubber resistor 164 and a snubber capacitor 168 are used to slow down the switch times. Also present is a built-in dead time to prevent shoot-through currents in the half-bridge. Though only one fluorescent lamp 140 is shown, additional fluorescent lamps may be incorporated in parallel with lamp 140, each with its own L-C resonant circuit. Any number of tube circuits can thus be driven from the single pair of power MOSFETs 20 and 21 if the power MOSFETs are sized accordingly to suit the resulting power level. The d-c bus also supplies power to the lamp failure protection arrangement of FIG. 3 which is comprised of a MOSFET 110 whose gate is coupled to a voltage divider formed by resistor 170, zener diode 108 and resistor 106 and having a source and drain that are connected across the V CC pin and the G pin of chip 30. According to this embodiment, when the lamp 140 is removed or fails, the voltage across the voltage divider rises and activates MOSFET 110 which clamps the V CC input to ground, thereby turning off the MOS gate driver chip 30 and cutting off the oscillating output supplied to the lamp 140. Alternatively, another type of MOS-gated device may be substituted for MOSFET 110. When the lamp 140 is replaced, the voltage across the voltage divider falls and turns off MOSFET 110. As a result, the operating voltage supplied to input V CC is reestablished so that chip 30 activates and again supplies power to MOSFETs 20 and 21 to ignite the lamp 140. FIG. 4 shows a ballast circuit which provides soft starting of two lamp ballast circuits and which also includes the novel lamp removal protection feature of the invention. Typically, the lamps are 32 W lamps. Before lamp 210 strikes, for example, the resonant circuit for the lamp comprises inductor 260 and capacitor 266, which are arranged in series with the lamp, as well as capacitor 242 which are arranged in parallel with the lamp. Similarly, the resonant circuit for lamp 212 includes inductor 260, capacitor 266 and capacitor 244. The ballast operating frequency is always above the resonance value. Initially, during pre-heat, the voltage across the lamps is too low for ignition. As the frequency drops, the lamp voltage increase until ignition occurs. At the same time, capacitor 234, which for example is a 220 microfarad capacitor and is arranged between the V CC and the COM terminals of the IC 30, charges from two current sources. The first current source is supplied at the terminal V+ by the diode bridge 202 and delivered through dropping resistor 222, which may be a 150 K-ohm resistor. The second current source is supplied by the ac output of the half-bridge MOSFETs 21 and 22 and delivered through a capacitor 237, which may be a 0.01 microfarad capacitor, that is arranged in series with a resistor 240, which is typically a 10 K-ohm resistor. As the capacitor 234 charges, V CC rises and reduces the output frequency of the half-bridge MOSFETS. This reduction in output frequency reduces the inductive reactance X L and raises the capacitive reactance X C so that the lamp voltage, which is dependent upon the ratio of X L and X C , rises with time until ignition occurs. The time required to pre-heat the lamp is thus a function of the charging time of capacitor 234. When one of the lamps 210 and 212 is removed or fails, the lamp circuit continues to operate. When both of the lamps 210 and 212 are removed or fail, the voltage across the voltage divider formed by resistor 230 and zener diode 252 increases and turns on MOSFET 250, thereby grounding the voltage V CC that is supplied to the input of the MOS gate driver chip 30. As a result, the gate voltages supplied to the gates of the half-bridge MOSFETs 20 and 21 fall and cut off the power supplied to the lamps. When one of the lamps is then replaced, the circuit will re-initiate the above-described soft start operation. FIG. 5 shows another example of a ballast circuit which incorporates another implementation of the soft starting and lamp removal protection features. Here, a compact fluorescent ballast circuit supplies power to a single lamp 310, which may be a 13 W to 26 W lamp. A driver IC 330, which for example is an International Rectifier IR51H420 integrated circuit that is suitable for a compact ballast circuit, incorporates the half-bridge connected MOSFETs that supply an ac output to the lamp at terminal V+. The remaining terminals of the IC operate in the same manner as the terminals having the same labels in the IC of the above embodiments. In this embodiment, the soft starting feature is provided by the resonant circuit that is formed by inductor 360 and capacitors 336 and 342 and that initially oscillates at a frequency higher than the operating frequency. The ratio of the L-C circuit elements causes pre-heating and ignition of the lamp 310 in the manner described above. Capacitor 334, which is arranged between the V CC and the COM terminals of the IC 330, is charged in a similar manner to that described in the circuit of FIG. 4 by current supplied by terminal V+ and diode bridge 302. The current is delivered through resistor 338 and causes V CC to rise until the lamp 310 ignites. When lamp 310 is removed, the voltage supplied to the gate of MOSFET 350 by the voltage divider formed by resistor 348, zener diode 346 and resistor 340 rises and activates the MOSFET. Thus, input voltage V CC is reduced to ground level and turns off driver integrated circuit 330 to protect the ballast circuit. Although the present invention has been described in relation to particular embodiments thereof, many other variations and modifications and other uses will become apparent to those skilled in the art. It is preferred, therefore, that the present invention be limited not by the specific disclosure herein, but only by the appended claims.

A lamp ballast circuit drives at least one gas discharge illumination device and includes first and second MOS-gated power semiconductor devices connected in a half-bridge configuration and coupled across a d-c supply with a common terminal located at the node between the devices for supplying an output signal to the illumination device. A self-oscillating driver circuit includes respective outputs for driving the power semiconductor devices and has at least one operating voltage supply terminal. When the illumination device fails or is removed, a protection circuit, electrically coupled to the illumination device, removes the operating voltage supplied to the supply terminal. The protection circuit includes a transformer/inductor coupled to the illumination device or a switch that is responsive to the voltage across the illumination device. Additionally, a soft-starting circuit can be provided for gradually increasing the voltage across the illumination device prior to igniting the illumination device.

CROSS REFERENCE TO RELATED APPLICATION [0001] This application claims priority from U.S. Provisional Application Ser. No. 61/658,736, filed Jun. 12, 2012, entitled “Recruiting Management System,” which is incorporated herein by reference. COPYRIGHT NOTICE [0002] This application includes material that is subject to copyright protection. The copyright owner has no objection to the facsimile reproduction by anyone of the patent disclosure, as it appears in the Patent and Trademark Office files or records, but otherwise reserves all copyright rights whatsoever. BACKGROUND [0003] One of the main challenges faced by a coach of a sports team (e.g., college sports team) lies in the recruiting of athletes. For example, colleges typically recruit athletes from junior colleges and/or high schools. Each year, the pool of available athletes is large, and the time that coaches or his/her staff can devote to recruiting activities is limited. Some athletes are more talented than others, and a coach must try to recruit the most talented athletes that he/she can get in the short amount of time available to him based on a limited scope of information. The information can be limited to an athlete's profile or information gathered at sporting events or recruiting events where the athlete's showcase their abilities (e.g., a scouting combine). [0004] The National Collegiate Athletic Association (NCAA) imposes recruiting process rules that in some manner inhibit a coach's ability to properly scout and accumulate all the proper and relevant information. Such rules make it difficult for a coach to establish a relationship with an athlete that the coach wants to recruit. Additionally, because statistics alone rarely tell the whole story, a coach usually will want to see an athlete in action before deciding whether to attempt to recruit that athlete. SUMMARY OF THE DISCLOSURE [0005] The present disclosure addresses failings in the art by providing a system and method for managing recruiting and tournament information gathered during the process of scouting. The present disclosure discusses an athlete recruiting system that allows a recruiting entity, such as a coach, to search for and/or identify athletes that satisfy recruiting criteria. Indeed, the disclosed methods and systems provide coaches the ability to adjust, update and view information related to prospects tangible attributes, event participation and performances, among other analytics associated with prospects. The information provided to a coach is communicated to a coach in an automated manner, which enables the coach to annotate, adjust and/or personally denote information personally observed by the coach in real-time. [0006] By way of background, from the perspective of both the recruiter (or coach) and the athlete, the current sports recruiting procedures suffer from several limitations. The required recruiting events necessarily place a geographical limitation on the entire process, restricting the pool of prospective athletes for particular teams and the choice of colleges for many high-school athletes. There is also a hierarchical limitation, smaller colleges with limited budgets and personnel are limited in the breadth of their search, and athletes from lesser known schools receive limited exposure to the recruiters. Additionally, the resources and tools provided to each coach at recruiting events are archaic and cumbersome to utilize. [0007] The recruiting system disclosed herein improves a way recruiters and coaches of sports teams can view, critique, and evaluate athletes at events. Recruiters, coaches, and, in some embodiments, athletes using the recruiting system (or portal) will be able to research and evaluate prospective athletes. Recruiters and coaches will have the ability to view an athlete's information including, but not limited to, pertinent statistics, biography, academic standing as well as, tangible and intangible attributes. Additionally, recruiters/coaches will have the ability to update and annotate an athlete's information based on the events. [0008] Additionally, using media streaming, recruiters/coaches will also be able to view associated video segments highlighting the athlete's participation in high school or college sports events in association with the event information accessed at said event. The recruiting system provides a portal that is also applicable to the transmission, storage and efficient access of information. Indeed, such information may also be accessible to the many different regulatory agencies involved in the governing of respective athletics. [0009] In accordance with one or more embodiments, a method and system are provided for recruitment management. In accordance with one or more embodiments, a non-transitory computer-readable storage medium is provided, the computer-readable storage medium tangibly storing thereon, or having tangibly encoded thereon, computer readable instructions that when executed cause at least one processor to perform steps within the scope of the present disclosure related to recruitment management. [0010] In accordance with one or more embodiments, a system is provided that comprises one or more computing devices configured to provide functionality in accordance with such embodiments. In accordance with one or more embodiments, functionality is embodied in steps of a method performed by at least one computing device. In accordance with one or more embodiments, program code to implement functionality in accordance with one or more such embodiments is embodied in, by and/or on a computer-readable medium. BRIEF DESCRIPTION OF THE DRAWINGS [0011] The foregoing and other objects, features, and advantages of the disclosure will be apparent from the following description of embodiments as illustrated in the accompanying drawings, in which reference characters refer to the same parts throughout the various views. The drawings are not necessarily to scale, emphasis instead being placed upon illustrating principles of the disclosure: [0012] FIG. 1 depicts an example of a system architecture according to some embodiments of the present disclosure; [0013] FIG. 2 depicts an example of a client device according to some embodiments of the present disclosure; [0014] FIG. 3 depicts a flow diagram according to some embodiments of the present disclosure; [0015] FIG. 4A-4E depict examples of embodiments displayed within a user interface of a computing device according to some embodiments of the present disclosure; [0016] FIG. 5 depicts an example of an embodiment according to the present disclosure; [0017] FIG. 6 depicts an embodiment displayed within a user interface of a computing device according to some embodiments of the present disclosure; and [0018] FIG. 7 is a block diagram illustrating an architecture of a hardware device in accordance with one or more embodiments of the present disclosure. DETAILED DESCRIPTION [0019] Subject matter will now be described more fully hereinafter with reference to the accompanying drawings, which form a part hereof, and which show, by way of illustration, specific example embodiments. Subject matter may, however, be embodied in a variety of different forms and, therefore, covered or claimed subject matter is intended to be construed as not being limited to any example embodiments set forth herein; example embodiments are provided merely to be illustrative. Likewise, a reasonably broad scope for claimed or covered subject matter is intended. Among other things, for example, subject matter may be embodied as methods, devices, components, or systems. Accordingly, embodiments may, for example, take the form of hardware, software, firmware or any combination thereof (other than software per se). The following detailed description is, therefore, not intended to be taken in a limiting sense. [0020] Throughout the specification and claims, terms may have nuanced meanings suggested or implied in context beyond an explicitly stated meaning. Likewise, the phrase “in one embodiment” as used herein does not necessarily refer to the same embodiment and the phrase “in another embodiment” as used herein does not necessarily refer to a different embodiment. It is intended, for example, that claimed subject matter include combinations of example embodiments in whole or in part. [0021] In general, terminology may be understood at least in part from usage in context. For example, terms, such as “and”, “or”, or “and/or,” as used herein may include a variety of meanings that may depend at least in part upon the context in which such terms are used. Typically, “or” if used to associate a list, such as A, B or C, is intended to mean A, B, and C, here used in the inclusive sense, as well as A, B or C, here used in the exclusive sense. In addition, the term “one or more” as used herein, depending at least in part upon context, may be used to describe any feature, structure, or characteristic in a singular sense or may be used to describe combinations of features, structures or characteristics in a plural sense. Similarly, terms, such as “a,” “an,” or “the,” again, may be understood to convey a singular usage or to convey a plural usage, depending at least in part upon context. In addition, the term “based on” may be understood as not necessarily intended to convey an exclusive set of factors and may, instead, allow for existence of additional factors not necessarily expressly described, again, depending at least in part on context. [0022] The present disclosure is described below with reference to block diagrams and operational illustrations of methods and devices to select and present media related to a specific topic. It is understood that each block of the block diagrams or operational illustrations, and combinations of blocks in the block diagrams or operational illustrations, can be implemented by means of analog or digital hardware and computer program instructions. These computer program instructions can be provided to a processor of a general purpose computer, special purpose computer, ASIC, or other programmable data processing apparatus, such that the instructions, which execute via the processor of the computer or other programmable data processing apparatus, implement the functions/acts specified in the block diagrams or operational block or blocks. In some alternate implementations, the functions/acts noted in the blocks can occur out of the order noted in the operational illustrations. For example, two blocks shown in succession can in fact be executed substantially concurrently or the blocks can sometimes be executed in the reverse order, depending upon the functionality/acts involved. [0023] For the purposes of this disclosure a computer readable medium (or computer-readable storage medium/media) stores computer data, which data can include computer program code (or computer-executable instructions) that is executable by a computer, in machine readable form. By way of example, and not limitation, a computer readable medium may comprise computer readable storage media, for tangible or fixed storage of data, or communication media for transient interpretation of code-containing signals. Computer readable storage media, as used herein, refers to physical or tangible storage (as opposed to signals) and includes without limitation volatile and non-volatile, removable and non-removable media implemented in any method or technology for the tangible storage of information such as computer-readable instructions, data structures, program modules or other data. Computer readable storage media includes, but is not limited to, RAM, ROM, EPROM, EEPROM, flash memory or other solid state memory technology, CD-ROM, DVD, or other optical storage, magnetic cassettes, magnetic tape, magnetic disk storage or other magnetic storage devices, or any other physical or material medium which can be used to tangibly store the desired information or data or instructions and which can be accessed by a computer or processor. [0024] The preferred embodiments of the present disclosure will now be described with reference to FIGS. 1-7 . FIG. 1 is a schematic diagram illustrating an example of an embodiment of a network. Other embodiments that may vary, for example, in terms of arrangement or in terms of type of components, are also intended to be included within claimed subject matter. As shown, FIG. 1 , for example, includes a variety of networks, such as local area local area network (LAN)/wide area network (WAN) 105 and wireless network 110 , a variety of devices, such as client device 101 and mobile device 102 , and a variety of servers, such as content server 107 and event server 106 . In connection with the illustrated example, which is non-exhaustive, additional or fewer servers may be utilized to transmit (and/or store data) between users, as will be understood from the below discussion. [0025] For purposes of this disclosure, a “server” should be understood to refer to a service point which provides processing, database, and communication facilities. By way of example, and not limitation, the term “server” can refer to a single, physical processor with associated communications and data storage and database facilities, or it can refer to a networked or clustered complex of processors and associated network and storage devices, as well as operating software and one or more database systems and application software that support the services provided by the server. [0026] In conjunction with the depiction illustrated in FIG. 1 , and discussed herein, servers may vary widely in configuration or capabilities, but generally a server may include one or more central processing units and memory. A server may also include one or more mass storage devices, one or more power supplies, one or more wired or wireless network interfaces, one or more input/output interfaces, or one or more operating systems, such as Windows Server, Mac OS X, Unix, Linux, FreeBSD, or the like. [0027] A content server may include a device that includes a configuration to provide content via a network to another device. A content server may, for example, host a site, such as a social networking site, examples of which may include, without limitation, Flicker, Twitter, Facebook, LinkedIn, or a personal user site (such as a blog, vlog, online dating site, etc.). A content server may also host a variety of other sites, including, but not limited to business sites, educational sites, dictionary sites, encyclopedia sites, wikis, financial sites, government sites, etc. [0028] A content server may further provide a variety of services that include, but are not limited to, web services, third-party services, audio services, video services, email services, instant messaging (IM) services, SMS services, MMS services, FTP services, voice over IP (VoIP) services, calendaring services, photo services, or the like. Examples of content may include text, images, audio, video, or the like, which may be processed in the form of physical signals, such as electrical signals, for example, or may be stored in memory, as physical states, for example. Examples of devices that may operate as a content server include desktop computers, multiprocessor systems, microprocessor-type or programmable consumer electronics, etc. [0029] A network may couple devices so that communications may be exchanged, such as between a server and a client device or other types of devices, including between wireless devices coupled via a wireless network, for example. A network may also include mass storage, such as network attached storage (NAS), a storage area network (SAN), or other forms of computer or machine readable media, for example. A network may include the Internet, one or more local area networks (LANs), one or more wide area networks (WANs), wire line type connections, wireless type connections, or any combination thereof. Likewise, sub networks, such as may employ differing architectures or may be compliant or compatible with differing protocols, may interoperate within a larger network. Various types of devices may, for example, be made available to provide an interoperable capability for differing architectures or protocols. As one illustrative example, a router may provide a link between otherwise separate and independent LANs. [0030] A communication link or channel may include, for example, analog telephone lines, such as a twisted wire pair, a coaxial cable, full or fractional digital lines including T1, T2, T3, or T4 type lines, Integrated Services Digital Networks (ISDNs), Digital Subscriber Lines (DSLs), wireless links including satellite links, or other communication links or channels, such as may be known to those skilled in the art. Furthermore, a computing device or other related electronic devices may be remotely coupled to a network, such as via a telephone line or link, for example. [0031] Although there are various types of networks, wireless networks may be used. A wireless may couple client devices with a network. A wireless network may employ standalone ad hoc networks, mesh networks, Wireless LAN (WLAN) networks, cellular networks, or the like. [0032] A wireless network may further include a system of terminals, gateways, routers, or the like coupled by wireless radio links, or the like, which may move freely, randomly or organize themselves arbitrarily, such that network topology may change, at times even rapidly. A wireless network may further employ a plurality of network access technologies, including Long Term Evolution (LTE), WLAN, Wireless Router (WR) mesh, or 2nd, 3rd, or 4th generation (2G, 3G, or 4G) cellular technology, or the like. Network access technologies may enable wide area coverage for devices, such as client devices with varying degrees of mobility, for example. [0033] For example, a network may enable RF or wireless type communication via one or more network access technologies, such as Global System for Mobile communication (GSM), Universal Mobile Telecommunications System (UMTS), General Packet Radio Services (GPRS), Enhanced Data GSM Environment (EDGE), 3GPP Long Term Evolution (LTE), LTE Advanced, Wideband Code Division Multiple Access (WCDMA), Bluetooth, 802.11b/g/n, or the like. A wireless network may include virtually any type of wireless communication mechanism by which signals may be communicated between devices, such as a client device or a computing device, between or within a network, or the like. [0034] Within the communications networks utilized or understood to be applicable to the present disclosure, such networks will employ various protocols that are used for communication over the network. Signal packets communicated via a network, such as a network of participating digital communication networks, may be compatible with or compliant with one or more protocols. Signaling formats or protocols employed may include, for example, TCP/IP, UDP, DECnet, NetBEUI, IPX, Appletalk, or the like. Versions of the Internet Protocol (IP) may include IPv4 or IPv6. The Internet refers to a decentralized global network of networks. The Internet includes local area networks (LANs), wide area networks (WANs), wireless networks, or long haul public networks that, for example, allow signal packets to be communicated between LANs. Signal packets may be communicated between nodes of a network, such as, for example, to one or more sites employing a local network address. A signal packet may, for example, be communicated over the Internet from a user site via an access node coupled to the Internet. Likewise, a signal packet may be forwarded via network nodes to a target site coupled to the network via a network access node, for example. A signal packet communicated via the Internet may, for example, be routed via a path of gateways, servers, etc. that may route the signal packet in accordance with a target address and availability of a network path to the target address. [0035] In some embodiments, the disclosed subject matter may comprise a content distribution network. A “content delivery network” or “content distribution network” (CDN) generally refers to a distributed content delivery system that comprises a collection of computers or computing devices linked by a network or networks. A CDN may employ software, systems, protocols or techniques to facilitate various services, such as storage, caching, communication of content, or streaming media or applications. Services may also make use of ancillary technologies including, but not limited to, “cloud computing,” distributed storage, DNS request handling, provisioning, signal monitoring and reporting, content targeting, personalization, or business intelligence. A CDN may also enable an entity to operate or manage another's site infrastructure, in whole or in part. [0036] Accordingly, in some embodiments, the present disclosure may be utilized via a content distribution system comprising a peer-to-peer network. A peer-to-peer (or P2P) network may employ computing power or bandwidth of network participants in contrast with a network that may employ dedicated devices, such as dedicated servers, for example; however, some networks may employ both as well as other approaches. A P2P network may typically be used for coupling nodes via an ad hoc arrangement or configuration. A peer-to-peer network may employ some nodes capable of operating as both a “client” and a “server.” [0037] According to some embodiments, the present disclosure may also be utilized within a social network. A social network refers generally to a network of individuals, such as acquaintances, friends, family, colleagues, or co-workers, coupled via a communications network or via a variety of sub-networks. Potentially, additional relationships may subsequently be formed as a result of social interaction via the communications network or sub-networks. A social network may be employed, for example, to identify additional connections for a variety of activities, including, but not limited to, dating, job networking, receiving or providing service referrals, content sharing, creating new associations, maintaining existing associations, identifying potential activity partners, performing or supporting commercial transactions, or the like. A social network may include individuals with similar experiences, opinions, education levels or backgrounds. Subgroups may exist or be created according to user profiles of individuals, for example, in which a subgroup member may belong to multiple subgroups. An individual may also have multiple associations within a social network, such as for family, college classmates, or co-workers. [0038] An individual's (e.g., coach's or athlete's) social network may refer to a set of direct personal relationships or a set of indirect personal relationships. A direct personal relationship refers to a relationship for an individual in which communications may be individual to individual, such as with family members, friends, colleagues, co-workers, or the like. An indirect personal relationship refers to a relationship that may be available to an individual with another individual although no form of individual to individual communication may have taken place, such as a friend of a friend, or the like. Different privileges or permissions may be associated with relationships in a social network. A social network also may generate relationships or connections with entities other than a person, such as companies, brands, or so-called ‘virtual persons.’ An individual's social network may be represented in a variety of forms, such as visually, electronically or functionally. For example, a “social graph” or “socio-gram” may represent an entity in a social network as a node and a relationship as an edge or a link. [0039] In some embodiments, multi-modal communications may occur between members of the social network. Individuals within one or more social networks may interact or communication with other members of a social network via a variety of devices. Multi-modal communication technologies refers to a set of technologies that permit interoperable communication across multiple devices or platforms, such as cell phones, smart phones, tablet computing devices, personal computers, televisions, SMS/MMS, email, instant messenger clients, forums, social networking sites, or the like. [0040] The above present persistent collaborative environment for interactive web applications provides persistence and sharing mechanisms for arbitrary application-defined objects. The sharing mechanism is coupled with a consistency mechanism that keeps client states consistent even when users perform conflicting operations. In addition, the framework maintains the interactivity of the web application at all times. The following sections provide descriptions of various embodiments of the present framework including the architecture, software and operations. [0041] FIG. 2 shows an example embodiment of a client (or user) device that may be used. A client device 200 may include a computing device capable of sending or receiving signals, such as via a wired or a wireless network. A client device may, for example, include a desktop computer or a portable device, such as a cellular telephone, a smart phone, a display pager, a radio frequency (RF) device, an infrared (IR) device, a web enabled Personal Digital Assistant (PDA), a handheld computer, a tablet computer, a laptop computer, a set top box, a wearable computer, a game console, smart TV, an integrated device combining various features, such as features of the forgoing devices, or the like. The user device (or client device) includes a processor and memory for storing and executing data and software. Computing devices may be provided with operating systems that allow the execution of software applications in order to manipulate data. A client device can be connected to the network, such as the Internet, via a wired data connection or wireless connection such as a Wi-Fi network, a satellite network or a cellular telephone network. A client device can support any type of interface for enabling the presentation or exchange of data. In addition, a user device may facilitate various input means for, but not limited to, receiving and generating information, including touch screen capability, keyboard and keypad data entry and voice-based input mechanisms. Any known and future implementations of user devices are applicable. [0042] The client device 200 may vary in terms of capabilities or features. Subject matter is intended to cover a wide range of potential variations. For example, a cell phone may include a numeric keypad or a display of limited functionality, such as a monochrome liquid crystal display (LCD) for displaying text. In contrast, however, as another example, a web-enabled client device may include one or more physical or virtual keyboards, mass storage, one or more accelerometers, one or more gyroscopes, global positioning system (GPS) or other location-identifying type capability, or a display with a high degree of functionality, such as a touch-sensitive color 2D or 3D display, for example. [0043] The client device may include or may execute a variety of operating systems, including a personal computer operating system, such as a Windows, iOS or Linux, or a mobile operating system, such as iOS, Android, or Windows Mobile, or the like. A client device may include or may execute a variety of possible applications, such as a client software application enabling communication with other devices, such as communicating one or more messages, such as via email, IMs, short message service (SMS), or multimedia message service (MMS), including via a network, such as a social network. A client device may also include or execute an application to communicate content, such as, for example, textual content, multimedia content, or the like. A client device may also include or execute an application to perform a variety of possible tasks, such as browsing, searching, playing various forms of content, including locally stored or streamed video, or games (such as fantasy sports leagues). The foregoing is provided to illustrate that claimed subject matter is intended to include a wide range of possible features or capabilities. [0044] A client device computing device may be capable of sending or receiving signals, such as via a wired or wireless network, or may be capable of processing or storing signals, such as in memory as physical memory states, and may, therefore, operate as a server. Thus, devices capable of operating as a server may include, as examples, dedicated rack-mounted servers, desktop computers, laptop computers, set top boxes, integrated devices combining various features, such as two or more features of the foregoing devices, or the like. [0045] By way of background, coaches main resource of information for recruiting are static and locally provided information sources. Generally, coaches can attend tournaments or other types of recruiting/scouting events to view prospects. As used herein, “prospects” refer to prospective athletes who coaches are scouting prospective of an athlete joining their team or being drafted by the coaches program or team. [0046] Generally, tournament or recruiting event organizers, referred to as “organizers” can create event. Such organizers can include coaches, athletic directors for universities or colleges, scouts for professional or semi-professional teams, high school coaches, AAU organizers or coaches, legion organizers or coaches, or any other type of person or entity involved in organizing a tournament or scouting event. Some non-limiting types of “events” include, but are not limited to, tournaments, scouting events, regular and/or post season tournaments, combines, and the like. For ease of explanation, reference will be made to “events” but should be understood to include any types of recruiting events that afford a potential scout, coach or “organizer” the ability to gauge and/or review an athletes ability (including in person meetings with the athlete and his/her family). Additionally, for each of explanation, recruiters, scouts or coaches will be referred to as “coaches”, in that, according to some exemplary embodiments, within the realm of collegiate athletics, coaches will most frequently be attending such events held by organizers. [0047] Currently, an organizer can organize or create an event (or tournament). The organizer, after organizing the event and sending out the necessary invitations to coaches and athletes, can create event data to that event. The event data file contains athlete and coach data, as discussed below. For example, the athlete data can include information related to each athlete's tangible and academic credentials. Coach data can comprise the program the coach is affiliated with, the program's ranking and affiliated division, the coach's record, and the coach's expressed needs (e.g., the positions or types of recruits the coach is interested in viewing). [0048] Generally, upon a coach attending an event, the coach will be given a binder of information related to the athletes present. The binder is typically a physical file of information related to the events ongoing and/or the athlete's present. That is, the information in the binder is a physical realization of the event data. Alternatively, the binder can comprise an electronic file (e.g., CD, disk, flash drive and the like) which a coach can load onto their device or print out. Such information is compiled by the organizer of the event, a third party, and is then physically handed over to the coach. As such, the coach must make manual notes in association with the binder, then proceed to process such information at a later time, after the event. Additionally, in order for the coaches to obtain the binder of event data the coach must pay a standard fee, in addition to the attendance fee (in some cases). [0049] The present disclosure remedies the shortcomings present in the field of recruiting at events, as discussed in reference to FIG. 3 . For example, after an organizer creates an event, the organizer can process the event data and store the data in a database which can be accessed locally or over a network. According to some exemplary embodiments, access to the event data is provided by a graphical user interface on a client's device via a downloadable application. The downloaded application, in accordance with the exemplary embodiments, renders a graphical user interface (GUI) which provides a user (e.g., coach, athlete or organizer) a visual display of event data, and the ability to access, view, render, modify, adjust, and/or update information, as discussed below and in relation to FIGS. 3-6 . For example, the recruiting system, which is provided/implemented via a downloaded application, can be downloaded from an “app store”, e.g., an online resource that provides downloadable application for user devices, or from a hosting web site (e.g., frontrush.com or coachpacket.com). For example, users (coaches) can download or use over-the-air applications or widgets (which are processes or functionality which run within an application, such as a client or server application) from a variety of online application vendors. Other embodiments may exist where the GUI for the recruiting system is provided by a client device, web-based provider and/or operating system running on the client device or server computing device. [0050] The event data, in a non-limiting example, can be stored in a the database(s) in the form of a CSV. In alternative embodiments, the data may be stored in an alternative format (e.g., XML, JSON, text, etc.) or maybe entered via a form (e.g., a web-based GUI). That is, the event data compiled by the organizer can be allocated according to a database of information associated with the attending athletes and coaches and uploaded to an event server, step 302 . It should be understood by those of ordinary skill in the art that the information comprised within the event data can be in the form of any readable/writable data that is importable and exportable from a database to and from a server(s), third party site, and a user's device. [0051] As discussed herein, the event data comprises information related to athletes. According to some embodiments, for each athlete, the event data includes specified types of information in a databases. Generally, the type of information stored is the kind of information that a recruiter or coach would be interested in knowing. For example, an athlete's information, some or all of which may comprise the athlete's profile, includes athletic performance information, athlete academic information, and athlete biographical information. Athlete performance information may include athletic statistical information and/or athletic performance video data. For example, an athlete's information may include the following types of information about the athlete: name, team position, home address, home telephone number, birth date, e-mail address, parent/guardian name, level (e.g., high school, junior college, etc.), school year (e.g., freshman, sophomore, junior, senior), years of NCAA eligibility, school name, school location (e.g., city and state), coach name, and coach telephone number. Furthermore, an athlete's information may include the following types of athlete academic information: grade point average and standardized test scores (e.g., ACT, SAT, PSAT, etc.). Additionally, an athlete's information may include the following types of athletic statistical information: height, weight, and sports-specific statistics (e.g., for basketball: points per game, rebounds per game, vertical jump height, field goal percentage, free throw percentage, 3-point field goal percentage, total assists, total steals, total blocks, etc.). An athlete's information also may include the following types of athlete biographical information: athletic achievements and awards, desired college qualities, intended major, desired sport program qualities, desired coach qualities, hobbies and activities, people who influence decisions (e.g., parents or guardians), and the like. According the some embodiments, the athlete themselves may provide the above information, or such information may be compiled over time during recruiting by a collection of coaches or a single coach or a collection of organizers or a single organizer, or any combination thereof. [0052] According to some embodiments, the event data may also comprise schedule information for the respective event or for attendees of the event. That is, for example, the event data may include information comprising, but not limited to, check-in/out times/dates for athletes and coaches, the schedule of games/activities taking place during the event, updated scores/results from such games/events, and/or any other type of data that can be compiled and communicated from such an event. In some embodiments, the event data, via the recruiting system, can indicate where or when during an event a coach or athlete is located. That is, the event data can provide a coach's or athlete's location within event and/or schedule for the event. The location of a coach/athlete can be visually displayed within the recruiting system's GUI on a map of the event. For example, Coach X is attempting to locate Athlete Y. Coach X identifies, via Athlete Y's event data (e.g., schedule information) that Athlete Y is playing on field 2 at the event at 11:00 am on Jun. 20, 2012. Coach X can also be presented with a layout/map of the event complex/facilities. The layout map can be provided on the GUI or within a portion of the GUI. On the map, an indicator providing Athlete Y's location can be depicted. The indicator can be annotated/or accompanied with appended or linked data providing the information or additional event data regarding Athlete Y. That is, the indicator may be clickable/selectable on the map, which can provide additional information about Athlete Y, for example, additional event data regarding the athlete's biography. [0053] The event data, according to some exemplary embodiments of the recruiting system discussed herein, can employ different encryption strengths and/or algorithms. In order to monetize this, organizers and/or content providers can structure different pricing/fees in accordance with encryption strength. For example, athlete, coach or program/event/organization-specific data, even when encrypted, can be stored in a separate directory from that of other athletes, coaches, or programs/events/organizations and/or users. In some embodiments, government rules and/or regulations can require different characteristics of encrypted data. The encryption can be a standard 256 bit AES (Advanced Encryption Standard) algorithm, approved by the NIST (National Institute of Standards and Technology), and uses both Symmetric and Asymmetric encryption/decryption keys. Thus, the communications associated with the recruiting system can be encrypted as well as any attachments or metadata (annotations, images, videos) associated with such communications. [0054] Upon an event's occurrence, the organizer can upload the event data as a file (a CSV as discussed above) to a central event server. Coaches, after registering with the event and receiving a user ID, can access the event data, step 304 . That is, according to some embodiments, coaches upon attending an event and/or in order to access said information will have register with the organizer or event. The coach can designate a specific identifier (e.g., the coach's email address) or be given an event specific identifier. [0055] Access to the event data can be based on the type and amount of information the coach desires. The event data can be accessed by the organizer denoting a specific network location to retrieve the event data. Alternatively, the organizer can distribute or alert attending coaches to a URL for which access to event data is given. In some embodiments, access to the event data can be based upon a coach's identifier and/or a tournament limit. That is, according to some exemplary embodiments, access to the event data is predicated upon a requesting coach paying an access fee, as seen in FIG. 5 . Further, as illustrated in FIGS. 4A-E , the event data can be displayed in on a graphical user interface (GUI) on a coach's device. In some embodiments, the display occurs within a browser running on the coach's device. [0056] For example, a coach can pay a certain amount for a single athlete's information, and/or a certain amount for a “bucket” (or plurality) of athletes. See FIG. 5 . For example, Coach A attends Event B and desires to download information related to Athletes X, Y and Z. As such, the organizer can charge the coach $0.99 per athlete, or a predetermined price for a “bucket”. That is, since there are 3 athletes listed in the coaches request, the coach can pay $2.50 for the bucket; a discounted rate due to the bucket. [0057] Accordingly, as discussed above, the coach can be given access to each athletes' information, steps. 306 and 308 . See also FIG. 4A . Access can be granted according to a centralized network location where access to each athletes' information is based upon the coach's identifier (e.g., username, email address, or personalized ID per event). That is, upon accessing the event, based on the coach's ID, the coach will only be afforded the information for which he/she has paid, step. 306 . In some embodiments, based on the information purchased by the coach, the organizer can provide the coach with a customized URL for which only the purchased or accessible information is present. [0058] According to some embodiments, the recruiting system can provide and display a listing of athletes that satisfy user-specified criteria, step 308 . For example, a coach can search for event data according to specified terms. This involves an input mechanism through which a user can input specified criteria. For example, the input mechanism may be an HTML form accessible via the recruiting system that includes multiple input fields. Each input field corresponds to an item of information that the athlete recruiting system stores in association with an athlete. For example, one input field may correspond to a “height” information item, and another input field may correspond to a “weight” information item. One or more input fields may correspond to sport-specific information. For example, an input field may correspond to a “vertical jump height” information item. Multiple input fields may correspond to a single information item. For example, both a “maximum height” input field and a “minimum height” input field may correspond to a single “height” information item. [0059] Accordingly, the HTML form includes a mechanism, such as a graphical button, that, when activated, causes a browser to send the contents of the input fields to an event server. In response to receiving the contents of the input fields, the event server generates and sends a database query to a database server. The database query instructs the database server to select, from a database, the identities of all of the athletes that satisfy the criteria specified in the input fields. The database query may also instruct the database server to select additional information from the available event data about the athletes that satisfy the criteria. In response to receiving the database query, the database server selects, from a database, the identities of all of the athletes that satisfy the criteria specified in the input fields. For example, if the criteria specify only a minimum height of 6′1″, then the database server selects all athletes associated with a height of at least 6′1″. Depending on the database query, the database server also may select, from the event data available on the database associated with the event server, additional information about the athletes that satisfy the criteria. The database server returns the selected information to the event server, which dynamically generates a web page that contains the selected information. The web server sends the dynamically generated web page to the browser from which the HTML form was submitted. The browser receives the web page and presents the web page to the coach that specified the criteria. Accordingly, access to such requested data can be predicated upon paying a fee for such data, either before searching or after performing said search. [0060] In one embodiment, the selected information includes the names of the athletes whose athlete information satisfies the specified criteria. Each such name may be presented in a separate link to a profile page for the corresponding athlete. See FIG. 4D . The links that indicate names of athletes whose profiles a recruiter or coach has already requested may be visibly distinguished from the links that indicate names of athletes whose profiles the recruiter has not already requested. For example, a browser may display, in a different color, links to pages that have already been visited. When generating the web page that lists the selected athlete names, an event server may use these associations to determine the manner in which the names of various athletes should be displayed, to distinguish athletes whose profiles have been requested in one or more of a particular organization's sessions from the names of athletes whose profiles have not been requested in any of the particular coaches sessions. According to one embodiment, the event server does not rank the selected information when generating the web page. This allows a recruiter to make his own judgments about the desirability of each athlete without being influenced by factors that might not be relevant to the recruiter. [0061] According to some embodiments, the recruiting system enables the coach to view the selected information and save and/or later restored said event data. For example, the web page indicating the results may contain a “save” control that, when activated, causes the information contained in the web page to be associated with a coach's session. Alternatively, the activation of the “save” control may cause the information contained in the web page to be stored in a file on the client device on which the browser displaying the web page resides. [0062] According to some embodiments, the recruiting system can enable users, such as coaches, organizers and/or athletes, to recognize other users using a augmented reality, virtual reality, and/or facial recognition, or a combination thereof, to immediately access event data. Augmented reality (AR), or mediated reality, is a live, direct or indirect, view of physical, real-world environment whose elements are augmented via computer-generated sensory input such as sound, video, images, graphics or GPS (location) data. The augmentation is typically in real-time and in semantic context with environmental elements, such as the event data, or athletes/coaches during the event. Thus, information related to an event, coach and/or athlete, e.g., event data, becomes interactive and provides a digitally manipulability environment. In some embodiments, event data can be overlaid on the GUI to provide a user with the full wealth of information available at an event and provided by the recruiting system discussed herein. [0063] For example, while a coach is attending an event and watching a game, a coach can be afforded the ability to point his device (e.g., mobile device) at an athlete and retrieve all the event data related to such athlete. That is, a coach can, for example, open the optics or infrared input/output of his/her device (e.g., image capture via a digital camera) and identify the identity of an athlete. As discussed above, the retrieved event data can be augmented on the coach's GUI. This enables the coach the ability to view information about the athlete on his/her device while watching the athlete. Augmentation can occur when a coach is viewing a video, image or other information of an athlete, and such information can be augmented with additional event data based on the image capture. [0064] By way of another non-limiting example, Coach X is watching a football game at an event and Athlete Y scores a touchdown. Coach X may want to identify the athlete and retrieve Athlete Y's information (or event data). Instead of searching the recruiting system, as discussed above, Coach X can point his/her device at Athlete Y and via facial recognition, Athlete Y's event data will be retrieved and visually displayed on Coach X's device. In some embodiments, the recognition of Athlete Y can be facial recognition. This can be predicated upon Athlete Y's event data including Athlete Y's picture/image, from which the image scanning of Athlete Y at the game can be compared against the images stored in the recruiting system database. It should be understood that recognition of an athlete should not be limited to facial recognition, in that an athlete or individual can be recognized based upon any traits which are recognizable from image capturing, or other known or to be know recognition/capturing techniques/functionality. In some embodiments, the scanning and retrieval of athlete data based on AR and/or recognition can be based upon proximity location and/or other information detection methods either known or to be known. In another example, the image capture of Athlete Y by Coach X can identify that Athlete Y is wearing a jersey with number 10 written on it. Based on the identification that number 10 has been identified in accordance with the location information derived from the image capture (e.g., one player wearing number 10 at field 2), or Coach X's location data (e.g., Coach X located at field 2), or the event data associated with an athlete wearing number 10 on a particular field (or location at the event—schedule information included within the event data), Athlete Y's data can be retrieved and displayed. Alternative embodiments also exist where event data can be retrieved based upon a scanned bar code (or other identifiable optical machine-readable representation of data) that may be worn on the athlete's jersey or within the proximity of the athlete (or location within/at the event). [0065] Additionally, as discussed above, coaches can access the event data by logging into an event via a specific identifier (e.g., the coach's email address) or be given an event specific identifier. Embodiments also exist where a coach can login to the event via facial recognition technology provided by the recruiting system, in that, upon the coach's image scan, the coach can be given access to the event data. In some embodiments, upon logging in, the coach can be displayed information related to event data regarding the event and/or athletes the coach has been scouting or viewing at the event, or in past events. [0066] In some embodiments, the recruiting system also allows athletes to search for organizations based on athlete-submitted criteria. Coaches may submit program (e.g., university/college) information to the recruiting system in the same manner that athlete information is compiled. The recruiting system may present an athlete with a list of programs or coaches that are associated with information that satisfies the athlete's criteria. Such criteria may include, for example, region, division, student body size, win-loss statistics, etc. [0067] As such, the coaches no longer need the physical binder of information. The information is accessible via the coaches device, e.g., tablet, smartphone and the like. That is, after registering with the event, receiving their ID (if needed) and paying for access to athlete information, coaches can immediately access event data based on the payment and authorization of the coach's ID. Indeed, this provides the coach with the ability view real-time information as the organizer can update the information about an event, in that, upon an update (e.g., an athlete's jersey number change), such updated data is immediately accessible to the viewing coach, step 310 . [0068] Additionally, in some embodiments, the recruiting system enables coaches, organizers and/or athletes (recruits) to communicate with each over via the recruiting system. That is, the recruiting system enables a coach, organizer or athlete to communicate with another respective user of the system via email, SMS, MMS, telephone, instant message, or via a social networking platform, such as, FACEBOOK®, TWITTER®, GOOGLE+®, and other known or to be known social and/or communication platforms. For example, after searching for an identifying an athlete that meets a coach's search criteria, or locating an athlete's event or biographical data within the recruiting system, the coach is afforded the functionality to communicate with the athlete via the recruiting system. Accordingly, the recruiting system enables such communication as discussed above, in that functionality for such communications are built in to the recruiting system's functionality (e.g., a widget, extension, ACTIVE X, or any other known or to be known enabling protocol). [0069] Indeed, in some embodiments, the recruiting system may only enable such communication in compliance with NCAA rules and regulations. That is, during a time period where a coach is allowed to contact a recruit, the recruiting system well allow such communications. However, should the coach attempt to communicate with an athlete during a restricted period as set by the NCAA, such communication will be denied. In such a situation, the coach will be alerted to the denial of contact with the athlete and in some embodiments, the coach can also be alerted the reasoning of such denial. That is, the recruiting system can provide the coach an alert that the communication is denied and provide the coach with information related to the NCAA rule governing such denial (in addition to a link to the supporting rules/regulations). In some embodiments, the NCAA can also be alerted to the attempted rule violation. Embodiments also exist involving the same as discussed above related to an athlete or organizer attempted to communicate with a coach. [0070] In some embodiments, as discussed herein, coaches can view the event data over a private network for the event or over a public network, e.g., the Internet. As such, the coaches will be afforded the ability to sync the event data they have accessed at the event with their other sources of information, step 312 . For example, a coach can annotate and make other notations about event data accessed at an event during the event (e.g., commenting on how an athlete is performing). See FIGS. 4A and 4B . Additionally, the coach can immediately sync said information with their account associated with their current recruiting software (e.g., applications provided by FrontRush™ at frontrush.com). In some embodiments, the recruiting system can allow events or organizers of events to market or advertise their events through a hosting site, e.g., frontrush.com or coachpacket.com. This enables coaches and athletes to search for and/or identify upcoming events. Also, the recruiting system can enable the coaches and athletes to indicate that they are attending the event, whereby such information can be communicated to the event organizer (or hosting entity). Additionally, in some embodiments, event data from previous events can be synchronized with event data at current events so that a coach, athlete and/or organizer is able to access and utilize an aggregate of information collected from a plurality of users using the recruiting system (from current or past events, or other platforms/resources where event type data can be collected). [0071] As such, the event data accessed by a coach can be stored in a “cloud” (which can be provided by the recruiting software provider (e.g., FrontRush™) or organizer). As used herein, a “cloud” is used in an art-recognized manner and can refer to a collection of centrally managed resources such as networked hardware and/or software systems and combinations thereof provided and maintained by an entity, wherein the collection of resources can be accessed by a user via wired or wireless access to a network that may be public or private, such as, for example, a global network such as the Internet. Such centralized management and provisioning of resources can provide for dynamic and on-demand provisioning of computing and/or storage to match the needs of a particular application. The cloud may include a plurality of servers, general or special purpose computers, as well as other hardware such as storage devices. The resources can include data storage services, word processing services, payment remitting services, and many other information technological services that are conventionally associated with personal computers or local and remote servers. Moreover, in one aspect, the resources can be maintained within any number of distributed servers and/or devices as discussed in more detail below. Thus, the present disclosure discusses a system that performs data management operations, including content-indexing and policy driven storage, within a cloud computing environment in order for a user to manage his/her personal information from a central online location. [0072] As discussed herein, the recruiting system can be a cloud-based data repository that functions as a secure off-site repository/data storage that, in some embodiments, synchronizes event data with other real-time or updated event data that can be accessed by a coach. Additionally, the event data can be synchronized with other information accessible to the coach (e.g., recruiting information associated with the coach on a recruiting site: FrontRush™ applications). Security for the recruiting system can be facilitated through the use of known cryptographic techniques, such as for example, symmetric and/or asymmetric cryptographic keying technology. [0073] Various monetization techniques or models may be used in connection with sponsored search advertising, including advertising associated with user search queries, or non-sponsored search advertising, including graphical or display advertising. In an auction-type online advertising marketplace, advertisers may bid in connection with placement of advertisements, although other factors may also be included in determining advertisement selection or ranking. Bids may be associated with amounts advertisers pay for certain specified occurrences, such as for placed or clicked-on advertisements, for example. Advertiser payment for online advertising may be divided between parties including one or more publishers or publisher networks, one or more marketplace facilitators or providers, or potentially among other parties. [0074] Some models may include guaranteed delivery advertising, in which advertisers may pay based at least in part on an agreement guaranteeing or providing some measure of assurance that the advertiser will receive a certain agreed upon amount of suitable advertising, or non-guaranteed delivery advertising, which may include individual serving opportunities or spot market(s), for example. In various models, advertisers may pay based at least in part on any of various metrics associated with advertisement delivery or performance, or associated with measurement or approximation of particular advertiser goal(s). For example, models may include, among other things, payment based at least in part on cost per impression or number of impressions, cost per click or number of clicks, cost per action for some specified action(s), cost per conversion or purchase, or cost based at least in part on some combination of metrics, which may include online or offline metrics, for example. [0075] A process of buying or selling online advertisements may involve a number of different entities, including advertisers, publishers, agencies, networks, or developers. To simplify this process, organization systems called “ad exchanges” may associate advertisers or publishers, such as via a platform to facilitate buying or selling of online advertisement inventory from multiple ad networks. “Ad networks” refers to aggregation of ad space supply from publishers, such as for provision en masse to advertisers. [0076] For web portals, advertisements may be displayed on web pages resulting from a user-defined search based at least in part upon one or more search terms. Advertising may be beneficial to users, advertisers or web portals if displayed advertisements are relevant to interests of one or more users. Thus, a variety of techniques have been developed to infer user interest, user intent or to subsequently target relevant advertising to users. One approach to presenting targeted advertisements includes employing demographic characteristics (e.g., age, income, sex, occupation, etc.) for predicting user behavior, such as by group. Advertisements may be presented to users in a targeted audience based at least in part upon predicted user behavior(s). [0077] Another approach includes profile-type ad targeting. In this approach, user profiles specific to a user may be generated to model user behavior, for example, by tracking a user's path through a web site or network of sites, and compiling a profile based at least in part on pages or advertisements ultimately delivered. A correlation may be identified, such as for user purchases or site visits, for example. An identified correlation may be used to target potential purchasers by targeting content or advertisements to particular users. Yet another approach to ad targeting could be based upon a user's current activity. Tracking of a user's path, or current location within a web site, or network of sites, and compiling of information associated with the user and/or the user's activity, or location. For example, regarding a user's activity, besides tracking which sites a user is viewing, information related to a user's activity within a site, e.g., their conversations with other users via the provided IM, provide advertisers with real-time topic points that warrant productive pivot points to provide real-time and relevant ads. [0078] An “ad server” comprises a server that stores online advertisements for presentation to users. “Ad serving” refers to methods used to place online advertisements on websites, in applications, or other places where users are more likely to see them, such as during an online session or during computing platform use, for example. [0079] During presentation of advertisements, a presentation system may collect descriptive content about types of content presented to users or the content being provided by the users on particular sites or via their interaction within a site/domain or network. A broad range of descriptive content may be gathered, including content specific to an advertising presentation system. Advertising analytics gathered may be transmitted to locations remote to an advertising presentation system for storage or for further evaluation. Where advertising analytics transmittal is not immediately available, gathered advertising analytics may be stored by an advertising presentation system until transmittal of those advertising analytics becomes available. [0080] Alternatively, the system may employ advertising techniques allowing individual advertisers to target users logging into a given session (or tournament) or accessing data from a provided resource or related to a single or plurality of athletes. In this embodiment, the system will allow an advertiser to monitor a coaches activity related to the information accessed or the location of such access. Also, such activity can be related to the type of information input or realized by each coach or user, or other users at the same event. In some embodiments advertisements may simply appear as a textual display optionally containing a hyperlink or other textual indicator. Alternatively, the advertisements may contain audio, video, or pictorial representations. In another embodiment, the system may allow advertisers to perform bulk analysis on past activity allowing for analytical analysis of the topics of discussion as a function of time. In another embodiment, advertisers may create interactive “bots” that interact with users (or coaches or prospects accessing the system) in manners other than simply broadcasting advertisements. In this example, an advertiser may utilize a bot that can simulate a discussion with the goal of advertising a particular event, service, product and/or assisting users in operating the website or webpage (e.g., coachpacket.com or other website hosting an event), as discussed in below. [0081] FIGS. 4A-4E depict non-limiting examples of embodiments of a user interface within the scope of the present disclosure. A browser executing on a client device displays a browser executable document (in some non-limiting examples a web page served from a website of a community provider system). The URL corresponding to the displayed web page is displayed by browser. As illustrated in FIG. 4A , a browser executable document is displayed to a coach on the coach's user device. In a non-limiting example according to FIG. 4A , an athlete's event data 402 is displayed to the coach on the coach's device. As discussed above, the coach can provide annotations 404 and 406 based on observations, which can be made based upon the athlete's performance at an event. As illustrated in FIG. 4B , the coach can add annotations 408 or comments about the athlete. As illustrated in FIG. 4C , the coach can access event data for past events 410 c , 410 d , 410 e the coach has attended. In some embodiments, upon payment of permission, the coach can access future event data 410 a , 410 b prior to events. As illustrated in FIG. 4D , the coach can access and view event data for a plurality of athletes 412 a - e , as discussed above. Accordingly, as illustrated in FIG. 4E , the coach can select particular athletes from the listing in FIG. 4D , and perform the steps illustrated in FIGS. 4A-4B , as discussed above. Accordingly, the depictions in FIGS. 4A-4E , in addition to the discussion above, can be performed on a mobile device, as illustrated in FIG. 6 . [0082] FIG. 7 is a block diagram illustrating an internal architecture of a computing device, e.g., a computing device such as server or user computing device, in accordance with one or more embodiments of the present disclosure. FIG. 7 illustrates a computer system upon which some exemplary embodiments of the present disclosure may be implemented. Although computer system 700 is depicted with respect to a particular device or equipment, it is contemplated that other devices or equipment (e.g., network elements, servers, processors) within can deploy the illustrated hardware and components of system 700 . [0083] As shown in FIG. 7 , internal architecture 700 includes one or more processing units, processors, or processing cores, (also referred to herein as CPUs) 712 , which interface with at least one computer bus 702 . Also interfacing with computer bus 702 are computer-readable medium, or media, 706 , network interface 714 , memory 704 , e.g., random access memory (RAM), run-time transient memory, read only memory (ROM), media disk drive interface 720 as an interface for a drive that can read and/or write to media including removable media such as floppy, CD-ROM, DVD, media, display interface 710 as interface for a monitor or other display device, keyboard interface 716 as interface for a keyboard, pointing device interface 718 as an interface for a mouse or other pointing device, and miscellaneous other interfaces not shown individually, such as parallel and serial port interfaces and a universal serial bus (USB) interface. [0084] Memory 704 interfaces with computer bus 702 so as to provide information stored in memory 704 to CPU 712 during execution of software programs such as an operating system, application programs, device drivers, and software modules that comprise program code, and/or computer-executable process steps, incorporating functionality described herein, e.g., one or more of process flows described herein. CPU 712 first loads computer-executable process steps from storage, e.g., memory 704 , computer-readable storage medium/media 706 , removable media drive, and/or other storage device. CPU 712 can then execute the stored process steps in order to execute the loaded computer-executable process steps. Stored data, e.g., data stored by a storage device, can be accessed by CPU 712 during the execution of computer-executable process steps. [0085] Persistent storage, e.g., medium/media 706 , can be used to store an operating system and one or more application programs. Persistent storage can also be used to store device drivers, such as one or more of a digital camera driver, monitor driver, printer driver, scanner driver, or other device drivers, web pages, content files, playlists and other files. Persistent storage can further include program modules and data files used to implement one or more embodiments of the present disclosure, e.g., listing selection module(s), targeting information collection module(s), and listing notification module(s), the functionality and use of which in the implementation of the present disclosure are discussed in detail herein. [0086] Network link 728 typically provides information communication using transmission media through one or more networks to other devices that use or process the information. For example, network link 728 may provide a connection through local network 724 to a host computer 726 or to equipment operated by an Network or Internet Service Provider (ISP) 730 . ISP equipment in turn provides data communication services through the public, world-wide packet-switching communication network of networks now commonly referred to as the Internet 732 . [0087] A computer called a server host 734 connected to the Internet 732 hosts a process that provides a service in response to information received over the Internet 732 . For example, server host 734 hosts a process that provides information representing video data for presentation at display 710 . It is contemplated that the components of system 700 can be deployed in various configurations within other computer systems, e.g., host and server. [0088] At least some embodiments of the present disclosure are related to the use of computer system 700 for implementing some or all of the techniques described herein. According to one embodiment, those techniques are performed by computer system 700 in response to processing unit 712 executing one or more sequences of one or more processor instructions contained in memory 704 . Such instructions, also called computer instructions, software and program code, may be read into memory 704 from another computer-readable medium 706 such as storage device or network link. Execution of the sequences of instructions contained in memory 704 causes processing unit 712 to perform one or more of the method steps described herein. In alternative embodiments, hardware, such as ASIC, may be used in place of or in combination with software. Thus, embodiments of the present disclosure are not limited to any specific combination of hardware and software, unless otherwise explicitly stated herein. [0089] The signals transmitted over network link and other networks through communications interface, carry information to and from computer system 700 . Computer system 700 can send and receive information, including program code, through the networks, among others, through network link and communications interface. In an example using the Internet, a server host transmits program code for a particular application, requested by a message sent from computer, through Internet, ISP equipment, local network and communications interface. The received code may be executed by processor 702 as it is received, or may be stored in memory 704 or in storage device or other non-volatile storage for later execution, or both. Other Embodiments [0090] In addition to the discussion above, the present disclosure embodies the following embodiments and functionality, as referenced from Appendix A of the instant application. According to some embodiments, tournament/recruiting event organizers can create an event and upload data (csv) to that event. This file contains athletes and athlete data. When they upload the file, they either pay us $0.99 an athlete or buy bucket i.e. x number of athletes in the file puts them in a specific bucket. Once the data is uploaded organizers or affiliates (third parties) can set a price (to charge coaches), choose whether coaches can export the data to csv, and/or limit which coaches can have access to the event (like a ‘guest list’). The event organizer is given a URL (or custom URL of our own) which they then distribute to coaches. The coaches can access this URL from their mobile device or from their laptop. If the tournament director charges the coaches, then the coaches can pay via credit card (or PayPal®) and the money transfers directly to the tournament organizer. The tournament director can also update data in real-time (e.g. an athlete's jersey number changes). [0091] At this point the coaches need online access to interact with the data but in the next few weeks, they will also have offline access. The coaches interaction includes searching for athletes, ranking athletes, and making notes on them. If they are also a Front Rush user, they can click ‘add to front rush’ and the data will transfer from their coach packet event to their Front Rush account. Any notes that were created on the recruit will show up in Front Rush as an “evaluation” as per NCAA requirements. The addition to Front Rush API is SOAP based or other known or to be known platform. Once the data is in Front Rush (e.g., existing or to be existing systems implemented, licensed, utilized or developed by FrontRush), it is used the same way any data in Front Rush is utilized. It also can be transferred to and from the schools respective admissions SIS (student information system). [0092] Additionally, the recruiting system provides additionally functionality, in accordance with some embodiments in accordance with the above disclosure: 1) Giving athletes the ability to login so they can keep their data current as well as pay to provide additional data that the coach will be able to see e.g. images, videos, Social data, other tournaments attended/attending. 2) Giving athletes the ability to take that login to other sites e.g. they use their “Coach Packet” login to fill out a recruit questionnaire that exists on a current schools site. 3) Giving athletes the ability to register and Pay for Coach Packet events with their login. Also, giving tournament directors the ability to manage the registration process. 4) Giving the athletes mobile access so they can update their events/data/and the like from the event. 5) Coach packet can be implemented, used as an extension of, and/or incorporated into the FrontRush™ or other recruiting tools or applications. 6) Native mobile versions of the recruiting system. [0093] For the purposes of this disclosure a module is a software, hardware, or firmware (or combinations thereof) system, process or functionality, or component thereof, that performs or facilitates the processes, features, and/or functions described herein (with or without human interaction or augmentation). A module can include sub-modules. Software components of a module may be stored on a computer readable medium for execution by a processor. Modules may be integral to one or more servers, or be loaded and executed by one or more servers. One or more modules may be grouped into an engine or an application. [0094] For the purposes of this disclosure the term “user”, “prospect”, “coach”, “student”, “athlete”, “subscriber” or “customer” should be understood to refer to a consumer of data supplied by a data provider. By way of example, and not limitation, the term “user” or “subscriber” can refer to a person who receives data provided by the data or service provider over the Internet in a browser session, or can refer to an automated software application which receives the data and stores or processes the data. [0095] Those skilled in the art will recognize that the methods and systems of the present disclosure may be implemented in many manners and as such are not to be limited by the foregoing exemplary embodiments and examples. In other words, functional elements being performed by single or multiple components, in various combinations of hardware and software or firmware, and individual functions, may be distributed among software applications at either the client level or server level or both. In this regard, any number of the features of the different embodiments described herein may be combined into single or multiple embodiments, and alternate embodiments having fewer than, or more than, all of the features described herein are possible. [0096] The disclosure described herein for recruiting athletes is applicable to a variety of contexts and implementations and is not limited to a particular context or implementation. For example, as used herein, organizations include recruiting entities such as colleges, universities, professional sports teams, semi-professional sports teams, apparel companies, scouting groups, marketing companies, sports agents groups, etc. Although approaches described herein are applicable to schools, approaches described herein are equally applicable to recruiting entities other than schools. Recruiting entities include entities that recruit athletes, but approaches described herein may be used to recruit individuals other than athletes. For example, approaches described herein may be used by business organizations to recruit employees. [0097] Functionality may also be, in whole or in part, distributed among multiple components, in manners now known or to become known. Thus, myriad software/hardware/firmware combinations are possible in achieving the functions, features, interfaces and preferences described herein. Moreover, the scope of the present disclosure covers conventionally known manners for carrying out the described features and functions and interfaces, as well as those variations and modifications that may be made to the hardware or software or firmware components described herein as would be understood by those skilled in the art now and hereafter. [0098] Furthermore, the embodiments of methods presented and described as flowcharts in this disclosure are provided by way of example in order to provide a more complete understanding of the technology. The disclosed methods are not limited to the operations and logical flow presented herein. Alternative embodiments are contemplated in which the order of the various operations is altered and in which sub-operations described as being part of a larger operation are performed independently. [0099] While various embodiments have been described for purposes of this disclosure, such embodiments should not be deemed to limit the teaching of this disclosure to those embodiments. Various changes and modifications may be made to the elements and operations described above to obtain a result that remains within the scope of the systems and processes described in this disclosure.

The present disclosure relates to a recruiting management system that acts as a computer network portal for recruiters, or recruiting coordinators, such as college and professional scouts/recruiters/coaches to disseminate, analyze and critique information related to prospects (e.g., high school, collegiate and/or professional athletes). In some embodiments, the recruiting management system enables recruiters and prospects real-time, up-to-date and on-line access to information related to recruiting events and/or statistical data. The management system facilitates recruiters the ability to communicate, search, select and screen the entire pool of athletes in a particular sport, based on customized and flexible criteria. The management system further enables recruiters the ability to adjust, update and view information related to prospects tangible attributes, event participation and performances, among other analytics associated with prospects.

DESCRIPTION OF THE INVENTION This application is a continuation-in-part of Ser. No. 285,937 filed July 23, 1981 which is a continuation of Ser. No. 142,154 filed Apr. 21, 1980 and is now abandoned. FIELD OF THE INVENTION This invention relates to container evacuation systems, and more particularly to a probe type of connector apparatus for use with flexible polymeric bag-type containers. DESCRIPTION OF THE PRIOR ART There has been an ever growing need for an inexpensive delivery system by which successive disposable containers of liquid food product can be connected to a delivery hose system and evacuated. The need for such a system has been greatest in the soft drink syrup industry, such as by fast food operators, bars, restaurants, and the like. In the past, soft drink bottlers have provided syrup to their customers in pressurized containers, typically in the form of metallic and plastic canisters. Such pressurized containers were then connected to the customer's liquid dispensing system. The liquid contents were then forced out of the containers and into the delivery tube system by a pressurized gas, typically carbon dioxide. Such prior art soft drink canisters, and the associated pressurized delivery system, had numerous disadvantages. One problem is that because these prior art canisters were typically formed from stainless steel, there were continual deterioration problems due to the fact that the highly corrosive syrup concentrations were in direct contact with the canisters' stainless steel walls. Another problem with such prior art pressurized canisters is that certain minimum pressure levels for the gas, such as carbon dioxide, are necessary to adequately force the soft drink product from the canister through the delivery tube system to the point of ultimate use. With certain diet soft drink syrups in which carbon dioxide is highly miscible, there oftentimes results in too much gas being entrained in the syrup due to the high gas pressure levels that are present. This results in poor taste characteristics for the finished soft drink product. Also, these pressurized canisters are oftentimes not entirely emptied in use, resulting in a continuous problem of residual product being left in the canisters and wasted. Further, use of such canisters is relatively expensive in that there are both high initial purchase costs involved as well as high transportation costs encountered in supplying canisters to the customer and returning them to the bottler. A more recent detrimental cost factor concerning such pressurized containers is the fact that the Federal Government has issued proposed guidelines under the Occupational Safety and Health Act which apparently labels them as "pressurized vessels," and as such, may require them to be annually inspected for safety reasons. Thus, the ability to use disposable flexible polymeric containers with liquid food product delivery systems has become important. However, up until the present invention, there have not been many satisfactory methods by which flexible bag containers could be effectively and inexpensively connected to a liquid product delivery system. (See. U.S. Pat. No. 4,137,930 for one known prior art method.) SUMMARY OF THE INVENTION These and other prior art problems have been overcome by the present invention. It provides a novel coupler apparetus having a displaceable seal plug type of fitment and a probe connector apparatus for use with flexible containers, such as foodstuff bags made of polymeric materials, and with associated liquid product delivery systems. This novel coupler apparatus utilizes both reusable components and disposable components. The disposable components include the flexible bag within which the product is contained and transported, a pouring nozzel or so-called fitment joined to the bag, and a cylindrical-shaped displaceable seal plug member which is slidably received within a passageway formed in the fitment. The reusable components are in the form of a probe connector permanently affixed at the connection end of the product delivery tube for a soft drink dispensing system. This reusable connector includes a probe member, a cylindrical sleeve and cylindrical inner cap comprising a probe adapter within which the probe is slidably retained and which is operable to detachably connect the probe connector to the fitment, and an outer cylindrical cap which is operable to move the probe member between opened and closed positions at the receiving end of the delivery tube, the outer cylindrical cap further being operable to move the seal plug member to open position with respect to the bag while opening the probe member with respect to the delivery tube. In use, the probe connector unit (with the probe member in its retracted position) is threadedly connected to the fitment of a flexible polymeric bag filled with soft drink syrup, for example. Once properly connected, the probe member is forced into the fitment by rotating the outer cylindrical cap thereby engaging and displacing the fitment's seal plug member farther into the fitment's passageway. This in turn exposes product drain means within both the seal plug member and probe member thereby allowing food product to flow from the bag into the delivery tube and on into the dispensing system. The liquid product can be delivered either by gravity flow or under the positive pressure of an associated pump. It is therefore a primary object of the present invention to provide a probe type of coupler apparatus for use with flexible food bags that are to be connected to liquid delivery systems, and which includes both reusable components as well as inexpensive disposable components. It is another object of the present invention to provide a fitment for a flexible foodstuff container which uses a displaceable seal plug and foil film seal to provide a tamper-proof product seal. It is yet another object to provide a fitment for a flexible polymeric container which has a displaceable seal plug member for eliminating the majority of the product from the fitment area thereby tending to reduce the overall oxygen transmission into the food product. It is a further object of the present invention to provide a probe type of coupler apparatus for soft drink syrup delivery systems having both leak-proof engagement to and dripless engagement from the fitment of a flexible container. It is a still further object to provide a probe type of coupler apparatus for a liquid dispensing system for a flexible foodstuff container whereby the probe is prohibited from inadvertently piercing a wall of the flexible container. It is yet a further object to provide a seal plug member for the fitment of a flexible foodstuff container which can not be inadvertently displaced into the interior of the container. The means by which the foregoing and other objects of the present invention are accomplished and the manner of their accomplishment will be readily understood from the following specification upon reference to the accompanying drawings, in which: FIG. 1 is a partially fragmented elevation view of a form of closure fitment member usable with the present invention; FIG. 2 is a form of seal plug member for use within the fitment shown in FIG. 1; FIG. 3 is a sectional elevation view of a flexible food bag fitted with the aforesaid closure fitment and seal plug members shown in a pre-fill position; FIG. 4 is a view similar to FIG. 3 showing the members in an intermediate, tamper-proof, post-fill position; FIG. 5 is a sectional elevation view with the probe connector in back-seated position and separated from the closure fitment member; FIG. 6 is a view similar to FIG. 5 showing the probe connector in an intermediate connected position; and FIG. 7 is a view similar to FIG. 6 with the parts shown in valve open, final drain position. DESCRIPTION OF THE PREFERRED EMBODIMENT Having reference to the drawings, wherein like reference numerals indicate corresponding elements, there is shown in FIG. 3 an illustration of a flexible polymeric food container bag 121. While not forming a part of the present invention, the bag is of the type having walls formed of multi-layered polymeric film (not shown) which typically are thermally bonded at their edges. (See U.S. Pat. Nos. 3,090,526; 3,556,816; and 4,085,244 for a detailed description of such flexible foodstuff containers.) A pouring nozzle of so-called closure fitment 126, best shown in FIG. 1, is inserted through an opening 127 formed in bag 121. The fitment includes a base flange portion 131 and a hollow cylindrical spout wall portion 133. The top side of base flange 131 is thermally bonded to the inner surface of the bag wall around opening 127. On the outside of the cylindrical spout wall portion 133 are formed a pair of axially spaced support rings 136a and 136b which are strengthened at 90° intervals by interposed radial ribs 137. When evacuating the bag 121, the wall 142 of a cardboard carton may be interposed between the base flange portion 131 and the lower support ring 136b (see FIGS. 5, 6 and 7). Additionally, during filling of bag 121, the lower ring 131b may support the fitment 126 between suitable yoke fingers such as are those designated 44 in FIG. 2 of parent application Ser. No. 285,937. As seen in FIG. 1, spout portion 133 has an inner annular rib 150. The upper and lower diagonal end surfaces of rib 150 respectively provide a stop shoulder 154 and a lock step shoulder 156, the purpose of both of which will be explained later. As shown in FIGS. 2 and 3, a seal plug member 162 has an upwardly open tubular body portion 164. An external cylindrical surface 250 is sealed by wiper ring portion 248 at the bottom of fitment wall 133. An upwardly open, coaxial inner tubular wall 165 is closed at the bottom by a cap portion 166. The top end surface 167 of wall 165 is engageable with the bottom surface 200 of a sleeve 202 as will be described. The inner cylindrical surface 204 of inner tubular wall 165 receives O-rings 206 and 208 as will be described. First drain means in the form of a plurality of drain holes 168 are formed through the annular portion 210. At the upper end of the tubular wall 164 an external flange 170 serves as a stop ring member. Seal plug member tubular wall 164 has an outer annular rib 172. Upper and lower diagonal end surfaces of rib 172 respectively provide a stop shoulder 174 and a lock step shoulder 176. Spaced below the rib 172 is a minor outer annular rib 211 having adjacent upper and lower diagonal end surfaces 213 and 216. The outer diameter 212 of tubular body portion 164 is appreciably smaller than the inner diameter 214 of spout portion 133 so the seal plug member 162 can move relatively freely within the closure fitment 126 except where movement is constrained by interference between the rib 150 and ribs 172 and 211. Groove 215 between ribs 172 and 211 is substantially the same size and shape as rib 150; thus, in the position of FIG. 3, rib 150 seals groove 215 and prevents entrance of outside contaminants into the bag through the clearance space 217. A cap liner 178 formed of pressure-adherent metal foil is placed across the open end of seal plug 162 to seal it temporarily. Seal plug member 162 is insertable different depths, to different operative positions, in cylindrical spout portion 133 as follows: (a) rib 150 sits in groove 215 in the pre-fill position of FIG. 3: (b) rib 172 is immediately below rib 150, with lock step shoulder 156 engaging stop shoulder 174, as shown in the intermediate, tamperproof position of FIG. 4; and (c) ribs 172 and 211 engage the lower portion of spout portion wall surface 214 in the final drain position shown in FIG. 7. Movement of seal plug member 162 into the bag 121 is limited by engagement of flange 170 with spout portion stop surface 173 (FIG. 7). FIGS. 5, 6 and 7 show a probe connector device generally designated 180. This comprises a probe member 182, a sleeve 202, an inner cap 184 and a manually rotatable outer cap 188. As will be seen, the sleeve 202 and inner cap 184 are made in two separate pieces for manufacturing and assembly convenience, but function as one piece to assemble the probe member 182 centrally within the fitment 126. Probe member 182 comprises a hollow tubular body portion 191 having a cylindrical head section 195 at one end. An end wall 192 closes the head section. There is a reduced diameter section 196 and a nipple 193 at the opposite end. A product delivery tube 100 is compressed onto the nipple by a ferrule 102. Additionally, the head section 195 has a second drain means, namely, a plurality of flow or drain holes 194 adjacent the closed end wall 192. O-rings 206 and 208 are seated in grooves 197 and 198 flanking the holes 194. In the final drain position of FIG. 7, the O-rings seal against the inner cylindrical wall 204 of inner tubular wall 165 on opposite sides of drain holes 168. Also, in the FIG. 7 position, a reduced diameter wiper ring portion 252 at the top of tubular wall 165 seals against the outer cylindrical wall 254 on head section 195. In the FIGS. 5 and 6 positions, these O-rings seal against the internal cylindrical wall surface 199 of sleeve 202. The sleeve 202 has a reduced diameter, cylindrical wall surface 181 slidably engaging the outer cylindrical surface 201 of probe member 182. An internal shoulder 203 on the sleeve engages an external shoulder 205 on the probe member in the FIGS. 5 and 6 positions. Inner cap member 184 is bell-shaped, having large and small diameter tubular sections 207 and 209 respectively, joined by an annular section 218. Section 209 has an inner cylindrical wall surface 220 closely fitted to an outer cylindrical wall surface 222 of the sleeve. Inner cap member 184 is fastened to sleeve member 202 for simultaneous longitudinal movement along probe member outer surface 201 as follows: Internal shoulder 224 at the left end of section 209 engages an external shoulder 226 on sleeve 202. At the opposite end, a retaining ring 228 seated in a groove in the sleeve engages the end of section 209. Thus, between shoulder 226 and retaining ring 228, inner cap 184 and sleeve 202 are held against relative movement and function for all practical purposes as a unit. Coarse Acme screw threads 183 and 185 connect inner cap 184 and sleeve 202 with fitment 126. Outer cap member 188 is bell-shaped, generally similar to the inner cap member 184 except larger. It has large and small diameter tubular sections 230 and 232 respectively, joined by an annular section 234. Section 230 is knurled as at 236 to facilitate rotating it manually. The inner wall 238 is smooth, cylindrical and slightly larger in diameter to move freely relative to the outer cylindrical surface 240 of inner cup member section 207. Section 232 has a central opening 244 in an end wall 242 within which probe member section 196 is rotatably journaled. Relative axial movement between members 188 and 182 is limited by shoulder 244 and retaining ring 246. Coarse Acme threads 187 and 189 threadedly connect inner and outer cap members 184 and 188. Thus, manual rotation of member 188 moves probe member 182 axially from the intermediate tamper-proof position shown in FIG. 4 to the fully opened final drain position shown in FIG. 7, and then back again to the back-seated position shown in FIG. 5 after bag 121 is emptied. Step-by-step use and operation will now be described. Step 1. Heat seal fitment 126 into the bag 121 and temporarily insert the seal plug member 162 in the spout portion 133, in the pre-fill position of FIG. 3. Foil 178 temporarily closes the open end of the seal plug member 162. Step 2. Prior to filling, remove seal plug member 162. Step 3. Mount the fitment 126 and bag 121 in any suitable support, for example, between fingers 44 of the fill support stand 46 shown in FIG. 2 of the patent application Ser. No. 285,937. Fill the bag. Step 4. Push seal plug member 162 into fitment 126 until their outer ends are flush. This is to the intermediate, tamper-proof, "post-fill" position of FIG. 4. In this position, the external annular rib 172 on the seal plug member has been forced inwardly past the inner annular rib 150 of the fitment spout portion 133. Stop shoulder 174 is retained beneath the lock step shoulder 156. The seal plug member cannot now be removed. Foil 178 is still in place. Step 5. Place the filled bag 121 in a carton having walls 142 (FIG. 6) and ship it to a restaurant, bar, or other use point. Step 6. At the use point, extend the fitment through a hole 128 in the carton and remove the foil seal 178. This position is shown at the left hand portion of FIG. 5 where the end cap portion 166 of seal plug member 126 is not yet fully extended into the bag 121. Further, the reusable probe connector 180 is in its back-seated position shown at the right hand portion of FIG. 5, still separated from the fitment 126. Any liquid remaining in probe member 182 and discharge tube 100 will be retained by O-rings 206 and 208. Step 7. Connect threads 183 of inner cap 184 to threads 185 of fitment 126. Rotate inner cap 184 until its end surface 186 is firmly, frictionally engaged with support ring 136a at the position shown in FIG. 6. Step 8. Grasp the knurled surface 236 of outer cap 188 and rotate it in a tightening direction. Cap 188 rotates and moves downwardly over the tubular section 207. This in turn moves the probe member 182 downwardly to the position shown in FIG. 7 with bottom end cap portion 166 and inlet ports 168 extended into the bag. Fluid can then drain from the bag to the delivery tube via holes 168 and 194, and probe member 182. Step 9. After the bag is emptied, rotate outer cap 188 in a loosening direction, upwardly to the FIG. 6 position. This is the "back-seated" (sealed) position referred to above in which liquid is positively retained within the probe 182. Step 10. Rotate inner cap 184 in a loosening direction, releasing the frictional engagement of its end surface 186 with the support ring 136a. Disconnect the reusable probe connector assembly 180 to the position shown in FIG. 5. Discard the empty bag and carton, fitment 126 and seal plug 162 and repeat this procedure with a new, filled bag 121. An important feature of the invention is the power screw assist provided by the threads 187 and 189 in moving the probe member 182 positively inwardly and outwardly between operative positions. Once bag 121 has been emptied, the probe member 182 can be back-seated into the sleeve 202 to positively close off the drain holes 194. Then, inner cap 184 can be unscrewed and the probe connector 180 disconnected for reuse. When the connector 180 is disconnected, as shown in FIG. 5, O-rings 206 and 208 are completely recessed within the end of the sleeve and protected from rough handling. As a further protection against abuse, the end 200 of the sleeve is itself recessed within the section 207 of the inner cap in the FIG. 5 position. From the foregoing, it is believed that those skilled in the art will readily appreciate the unique features and advantages of the present invention over previous types of fitments and couplers for flexible foodstuff bags. Further, it is to be understood that while the present invention has been described and illustrated with a particular preferred embodiment, as set forth in the accompanying drawings and as above described, the same nevertheless is susceptible to change, variation and substitution of equivalents without departing from the spirt and scope of this invention which should not be restricted by the foregoing description and drawings except as may appear in the following appended claims.

A connector apparatus for use with various types of disposable flexible foodstuff containers, such as polymeric bags containing soft drink syrups, of the type having a fitment attached to a wall of the container. The novel fitment includes a displaceable valve means having a seal plug member slidably received within a passageway and capable of being releasably retained at a pre-fill position and lockably retained at a post-fill position. A mating probe connector is permanently attached to the food product delivery hose system and includes an inner cap and sleeve comprising a probe adapter capable of being detachably secured to the novel fitment and a probe member for engaging the seal plug member. Product drain and positive evacuation structures are provided to assure substantially complete draining of product from the container during use. The probe is movable in response to rotation of a screw-threaded outer cap member.

FIELD OF THE INVENTION The present invention encompasses radiolabeled peptide analogs of vasoactive intestinal peptide (VIP) labeled with a radionuclide useful for imaging target sites within mammalian living systems. The invention particularly provides radiolabeled VIP derivatives that bind selectively to the VIP receptor on target cells. Specifically, the invention relates to the radiolabeling of VIP-receptor specific agents and their subsequent use for radiodiagnostic and radiotherapeutic purposes. The invention encompasses methods for radiolabeling these peptides with radionuclides and the use of these peptides as scintigraphic imaging agents. The radiolabeled VIP derivatives of the present invention exhibit pharmacological activity and therefore are useful as either imaging agent for visualization of VIP-receptor positive tumors and metastases, as a radiodiagnostic agent or as a radio-therapeutic agent for the treatment of such tumors in vivo by specifically targeting the cytotoxic radionuclide selectively to the tumor site in mammalian living systems. BACKGROUND OF THE INVENTION Vasoactive intestinal peptide is a 28-amino acid neuropeptide, which was first isolated from the porcine intestine (Said and Mutt, 1970). It bears extensive homology to secretin, PHI and glucagon. The amino acid sequence for VIP is: His-Ser-Asp-Ala-Val-Phe-Thr-Asp-Asn-Tyr-Thr-Arg-Leu-Arg-Lys-Gln-Met-Ala-Val-Lys-Lys-Tyr-Leu-Asn-Ser-Ile-Leu-Asn-NH 2 (SEQ ID NO: 1) VIP is known to exhibit a wide variety of biological activities such as the autocrine, endocrine and paracrine functions in living organisms (Said, 1984). In the gastrointestinal tract, it has been known to stimulate pancreatic and biliary secretions, hepatic glycogenesis as well as the secretion of insulin and glucagon (Kerrins and Said, 1972; Domschke et al., 1977). In the nervous system it acts as a neurotransmitter and neuromodulator, regulating the release and secretion of several key hormones (Said, 1984). In recent years, attention has been focussed on the function of VIP in certain areas of the CNS as well its role in the progression and control of neoplastic disease (Reubi, 1995). The importance of peptide growth factors and regulatory hormones in the etiology and pathogenesis in several carcinomas has long been recognized. Data from epidemiological and endocrinological studies suggest that neuropeptides like VIP which are responsible for the normal growth of tissues like the pancreas can also cause conditions for their neoplastic transformation (Sporn et al., 1980). Several lines of evidence indicate that VIP acts as a growth factor and plays a dominant autocrine and paracrine role in the sustained proliferation of cancer cells (Said, 1984). The stimulatory effect of VIP on tumor growth can be mediated directly by its receptors on cell membranes or indirectly by potentiation of the activities of other growth factors in tumor cells (Scholar E. M. Cancer 67(6): 1561-1569, 1991). The synergistic effect of VIP and related pituitary adenylate cyclase activating polypeptide (PACAP) in glioblastomas is an illustration to the above fact (Moody, T. W., et al. Peptides 17(3), 545-555, 1996). The multiple physiological and pharmacological activities of VIP are mediated by high affinity G-protein coupled transmembrane receptors on target cells (Reubi, 1995). VIP receptors are coupled to cellular effector systems via adenylyl cyclase activity (Xia et al., 1996). The VIP receptor, found to be highly over-expressed in neoplastic cells, is thought to be one of the biomarkers in human cancers (Reubi, 1995). High affinity VIP receptors have been localized and characterized in neoplastic cells of most breast carcinomas, breast and prostate cancer metastases, ovarian, colonic and pancreatic adenocarcinomas, endometrial and squamous cell carcinomas, non small cell lung cancer, lymphomas, glioblastomas, astrocytomas, meningiomas and tumors of mesenchymal origin. Amongst, neuroendocrine tumors all differentiated and non-diffemtiated gastroenteropancreatic tumors, pheochromocytomas, small-cell lung cancers, neuroblastomas, pituitary adenomas as well tumors associated with hypersecretory states like Verner-Morrison syndrome were found to overexpress receptors for vasoactive intestinal peptide (Reubi, 1995, 1996, 1999; Tang et al., 1997a&b; Moody et al., 1998a&b; Waschek et al., 1995; Oka et al., 1998)). These findings suggest that new approaches for the diagnosis and treatment of these cancers may be based on functional manipulation of VIP activity, using synthetic peptide analogs of the same. Historically, the somatostatin analog 111 In-DTPA-[D-Phe 1 ]-octreotide is the only radiopeptide, which has obtained regulatory approval in USA and Europe (Lamberts et al., 1995). Radiolabeled VIP has been shown to visualize a majority of gastropancreatic adenocarcinomas, neuroendocrine tumors, as well as insulinomas (which are often missed by radiolabeled octreotide) (Behr et al., 1999). VIP-receptor scinitigraphy offers certain advantages over radioimaging involving somatostatin receptors. The presence of high affinity receptors for VIP have been demonstrated in a larger number of human tumors, relative to the somatostatin receptors. Secondly, the density of VIP receptors on tumors has been found to be greater than somatostatin (Behr et al., 1999). Therefore, the VIP-receptor scan is more sensitive and convenient in localizing tumors and their metastatic spread as compared to somatostatin. The applications of this technique are manifold. It has been used for the sensitive detection of VIP-receptor positive tumors. This includes primary carcinoids, cancers of the gastrointestinal tract as well as distant metastases (Reubi, 1995, 1996). It can also be used to target cytotoxic radionuclides specifically to the tumor site. It predicts the VIP-receptor status of the patient and thereby the response of the patient towards radiotherapy by radiolabeled VIP analogs. Lastly, such radiolabeled peptides have been successfully used in radioguided surgery (Lamberts et al., 1995). 123 I-VIP, 125 I-VIP and their derivatives have been extensively used for imaging pancreatic adenocarcinomas, endocrine tumors of the gastrointestinal origin, mesenchymal tumors as well secondary tumor metastatic sites, in patients (Jiang et al., 1997; Virgolini et al., 1996, 1998; Raderer et al., 1998 ; Moody et al., 1998; Kurtaran et al., 1997; Pallella et al., 1999). Radioiodinated VIP and its derivatives have been also used to assess the binding affnty of peptides for VIP-receptors on tumor cells in vitro. The biodistribution, safety and absorbed dose of the aforesaid radioiodinated peptide derivatives have also been studied earlier (Virgolini et al., 1995). U.S. Pat. No. 5,849,261, granted to Dean et al., on Dec. 15, 1998 describes the applications of radiolabeled vasoactive intestinal peptide (VIP) for diagnosis and therapy. In particular, this U.S. Patent discloses a method for preparing a radiopharmaceutical agent, comprising native vasoactive intestinal (VIP) peptide attached to a radionuclide like technetium or rhenium via a chelating moiety. The radiopharmaceutical when labeled with technetium or rhenium via a chelating moiety has a VIP binding affinity which is not less than about one tenth the affinity of radioiodinated native VIP for the receptor. However, there is still a need for improved synthetic analogs of VIP as radiopharmaceuticals, which are easy to generate and are capable of being employed with higher sensitivity and specificity in terms of their radioimaging and radiodiagnostic properties. This invention describes the preparation and use of peptide analogs of VIP having constrained amino acids. The design of conformationally constrained bioactive peptide derivatives has been one of the widely used approaches for the development of.peptide-based therapeutic agents. Non-standard amino acids with strong conformational preferences may be used to direct the course of polypeptide chain folding, by imposing local stereochemical constraints, in de novo approaches to peptide design. The conformational characteristics of α,α-dialkylated amino acids have been well studied. The incorporation of these amino acids restricts the rotation of φ, Ψ angles, within the molecule, thereby stabilizing a desired peptide conformation. The prototypic member of α,α-dialkylated aminoacids, α-amino-isobutyric acid (Aib) or α,α-dimethylglycine has been shown to induce β-turn or helical conformation when incorporated in a peptide sequence (Prasad and Balaram, 1984, Karle and Balaram, 1990). The conformational properties of the higher homologs of α,α-dialkylated amino acids such as di-ethylglycine (Deg), di-n-propylglycine (Dpg), di-n-butylglycine (Dbg) as well as the cyclic side chain analogs of α,α-dialkylated amino acids such as 1-aminocyclopentane carboxylic acid (Ac5c), 1 -aminocyclohexane carboxylic acid (Ac6c), 1-aminocycloheptane carboxylic acid (Ac7c) and 1-aminocyclooctane carboxylic acid (Ac8c) have also been shown to induce folded conformation (Prasad et al., 1995 ; Karle et al., 1995). α,α-dialkylated amino acids have been used in the design of highly potent chemotactic peptide analogs (Prasad et al., 1996). The present invention incorporates the conformational properties of such α,α-dialkylated amino acids for the design of biologically active peptide derivatives, taking VIP as the model system under consideration. REFERENCES Behr T. M. et al. Q. J. Nucl. Med., 43, 268-280,1999. Domschke, S. et al. Gastroenterology, 73, 478-480, 1977. Jiang S. et al. Cancer Res., 57, 1475-1480,1997. Karle, L L. et al. (1995) J. Amen Chem. Soc. 117, 9632-9637. Karle, L L. and Balaram, P. (1990) Biochemistry 29, 6747-6756. Kerrins, C. and Said, S. I. Proc. Soc. Exp. Biol. Med., 142, 1014-1017, 1972. Kurtaran A. et al. J. Nucl. Med., 38, 880-881, 1997. Lamberts, S. W. J. et al. In Somatostatin and its Receptors, Ciba Found. Symp., 190, 222-239, 1995. Moody, T. W. et al. Peptides, 19 (3), 1998a. Moody, T. W. et al. Ann. N.Y. Acad. Sci., 865, 290-296. 1998b. Oka, H. et al. Am. J. Pathol., 153 (6), 1787-1796, 1998. Pallella, V. R. et al. J. Nucl. Med., 40(2), 352-360, 1999. Prasad, B. V. V and Balaram, P. (1984) CRC Crit. Rev. Biochem. 16, 307-347. Prasad, S et al. (1995) Biopolymers 35, 11-20. Prasad, S et al. (1996) Int. J. Peptide Protein Res. 48, 312-318. Reubi, J. C. J. Nucl. Med., 36 (10), 1995. Reubi, J. C. et al. Int. J. Cancer, 81 (3), 1999. Reubi, J. C. et al. Cancer Res., 56 (8), 1922-1931, 1996. Raderer, M. et al. J. Nucl. Med., 39 (9), 1570-1575, 1998. Said, S. I. and Mutt, V. Science, 169, 1217-1218,1970. Said, S. I. Peptides, 5, 143-150, 1984. Sporn, M. B., and Todaro, G. J. N. Engl. J. Med., 303, 378-379, 1980. Tang, C. et al., Gut, 40 (2), 267-271, 1997a. Tang, C. et al., Br. J. Cancer, 75 (10)1467-1473, 1997b. Virgolini, I. et al. J. Nucl. Med., 36(10), 1732-1739, 1995. Virgolini, I. et al. Nucl. Med. Biol., 23 (6), 685-692, 1996. Virgolini, I. et al. J. Nucl. Med., 39 (9), 1998. Waschek, J. A. et al. Cancer Lett., 92 (2), 1995. Xia, M. et al., J. Clin. Immunol., 16 (1), 21-30, 1996. Throughout the specification and claims the following abbreviations are used: Aib: α-Aminoisobutyric acid Deg: α,α-Diethylglycine Ac5c: 1-Amino Cyclopentane Carboxylic acid BOP: Benzotriazole-1-yl-oxy-tris-(dimethylamino)-phosphonium hexofluorophosphate PyBOP: Benzotriazole-1-yl-oxy-tris-pyrrolidino-phosphonium hexofluorophospate HBTU: O-Benzotriazole-N,N,N′,N′-tetramethyl-uronium-hexofluoro-phosphate TBTU: 2-(1H-Benzotriazole-1-yl)-1, 1, 3, 3-tetramethyluronium tetrafluoroborate HOBt: 1-Hydroxy Benzotriazole DCC: Dicyclohexyl carbodiimide DIPCDI: Diisopropyl carbodiimide DIEA: Diisopropyl ethylamine DMF: Dimethyl formamide DCM: Dichloromethane NMP: N-Methyl-2-pyrrolidinone TFA: trifluoroacetic acid Throughout the specification and claims the amino acid residues are designated by their standard abbreviations. Amino acids denote L-configuration unless indicated by D or DL appearing before the symbol and separated from it by a hypen. SUMMARY OF THE INVENTION The present invention encompasses radiolabeled peptide analogs of vasoactive intestinal peptide (VIP) labeled with a radionuclide useful for imaging target sites (e.g. use as a scintigraphic imaging agent), for use in radiodiagnostics, and radiotherapy within mammalian living systems. The invention particularly provides radio labeled VIP derivatives that bind selectively to the VIP receptor on target cells. Specifically, the invention relates to the radiolabeling of VIP receptor specific agents and their subsequent use for radiodiagnostic and radiotherapeutic purposes. The invention encompasses methods for radiolabeling these peptides with radionuclides and the use of these peptides as scintigraphic imaging agents. The present invention also encompasses the use of these radiolabeled peptides as anti-neoplastic agents for specific radiotherapy in cancer. A further object of the invention is the use of certain novel VIP analogs to determine the binding affinities of these peptides for their cognate receptors on cancer cells. BRIEF DESCRIPTION OF THE FIGURE FIG. 1 shows a radioimage of a mammalian tumor xenografted nude mouse using a radiolabeled VIP analog. DETAILED DESCRIPTION OF THE INVENTION The present invention provides novel radiolabeled peptide analogs of vasoactive intestinal peptide useful for imaging target sites within a mamalian living system, comprising a synthetic receptor-binding peptide analog of vasoactive intestinal peptide (VIP) radiolabeled with a radionuclide. The present invention relates to the use of the radiolabeled peptides, processes for production of the radiolabeled peptides, pharmaceutical preparations for its use as a diagnostic, imaging as well as a radiotherapeutic agent in vivo. The VIP peptide analogs of the present invention, which is a VIP receptor antagonist has the sequence: His-Ser-Asp-Xxx-Val-4-Cl-D-Phe-Thr-Asp-Asn-Tyr- Thr-Arg-Leu-Arg-Lys-Gln-Leu-Ala-Val-Lys-Lys-Tyr- Leu-Asn-Ser-Ile-Leu-Asn-NH 2 where Xxx is Aib, Deg or Ac5c or a pharmaceutically acceptable salt of the peptide. A preferred peptide is: His-Ser-Asp-Aib-Val-4-Cl-D-Phe-Thr-Asp-Asn-Tyr- Thr-Arg-Leu-Arg-Lys-Gln-Leu-Ala-Val-Lys-Lys-Tyr- Leu-Asn-Ser-Ile-Leu-Asn-NH 2 (SEQ ID NO: 2) or a pharmaceutically acceptable salt thereof. Salts encompassed within the term “pharmaceutically acceptable salts” refer to non-toxic salts of the compounds of this invention. Representative salts and esters include: acetate, ascorbate, benzenesulfonate, benzoate, bicarbonate, bisulfate, bitartrate, borate, camsylate, carbonate, citrate, dihydrochloride, methanesulfonate, ethanesulfonate, p-toluenesulfonate, cyclohexylsulfamate, quinate, edetate, edisylate, estolate, esylate, fumaxate, gluconate, glutamate, glycerophophates, hydrobromide, 5 hydrochloride, hydroxynaphthoate, lactate, lactobionate, laurate, malate, maleate, mandelate, mesylate, mucate, napsylate, nitrate, n-methylglucamine, oleate, oxalate, palmoates, pamoate (embonate), palmitate, pantothenate, perchlorates, phosphate/diphosphate, polygalacturonate, salicylates, stearate, succinates, sulfate, sulfamate, subacetate, succinate, tannate, tartrate, tosylate, trifluoroacetate and valerate. Other salts include Ca, Li, Mg, Na and K salts; salts of amino acids such lysine or arginine; guanidine, diethanolamine or choline; ammonium, substituted ammonium salts or aluminum salts. The salts can be prepared by standard techniques. The VIP receptor antagonist: His-Ser-Asp-Xxx-Val4-Cl-D-Phe-Thr-Asp-Asn-Tyr-Thr-Arg-Leu-Arg-Lys-Gln-Leu-Ala-Val-Lys-Lys-Tyr-Leu-Asn-Ser-Ile-Leu-Asn-NH 2 where Xxx is Aib, Deg or Ac5c have been shown in co-pending application Ser. No. 09/630,335 (filed on Jul. 31, 2000) to be selectively binding to VIP receptors on cancer cells. The anti-proliferative activity of the aforesaid VIP antagonist has been previously demonstrated in a number of experimental models of pancreatic, prostate, mammary and lung cancer, suggesting its high anti-neoplastic therapeutic potential. The applicants have found that the VIP analogs of the present invention have greater affinity for its cognate receptors on tumor cells as compared to native VIP, which in turn leads to better radio-imaging, radiodiagnostic and radio-therapeutic efficacy of the radiopharmaceuticals of the present invention. While not wishing to be bound by theory, the applicants believe that the improved efficacy of the radiopharmaceuticals of the present invention are due to the nature of the VIP analogs themselves, which have receptor bound conformations caused by the incorporation of the unusual amino acids. The labeling of peptides by a radionuclide has been accomplished in the present invention, by several strategies: 1. Direct labeling of radionuclide to the peptide analogs. 2. attachment of chelating groups to the peptide and subsequent radiolabeling by radionuclide. 3. Incorporation of radionuclide to chelator moieties covalently linked to the peptide via a spacer group. It is important to note that in the above cases, the chelator and spacer groups are incorporated site-specifically at a position which does not affect the specific binding properties of the peptide to the VIP receptor on tumor cells in vitro and in vivo. In a preferred embodiment of the present invention, the radionuclide is selected from Technetium (Tc-99m), Iodine 123 ( 123 I), Iodine 131 ( 131 I), Indium-111 ( 111 I) and Rhenium-188 ( 188 Re). One embodiment of the invention involves the radiolabeling of the VIP antagonists directly by a radionuclide such as Tc-99m. Tc-99m forms a coordinate covalent linkage with certain specific amino acid residues of the peptide. The formation of a stable Tc-peptide bond is one of the major advantages for its use for imaging purposes. The attachment of Tc-99m to the peptide involves the reaction of a salt of Tc-99m such as pertechnate to the peptide, in the presence of a reducing agent such as dithionate ion, ferrous ion or stannous chloride. The radiolabeled peptides are separated from the unincorporated Tc-99m as described in the examples and used for radioscintography. Another embodiment of the present invention includes the attachment of certain chelating groups to the VIP analog. The chelating groups are capable of complexing a detectable element such as a radionuclide. According to the invention, the chelating moiety may be attached directly or indirectly to the peptide, e.g. by means of a spacer or a bridging group to the amino terminus of the VIP analog. In a more preferred embodiment of the invention the radionuclide is Tc-99m bound to a chelating moiety. All the radiolabeled chelated peptides retain their affinity for VIP receptors on cancer cells. According to one embodiment of the invention, the chelating group has substantial hydrophobic character. Examples of chelating groups include e.g. iminodicarboxylic groups, polyarninopolycarboxylic groups, e.g. those derived from non-cyclic ligands e.g. ethylene diamine tetra acetic acid (EDTA), diethylene triamine pentaacetic acid (DTPA), ethylene glycol-0,0′-bis (2-aminoethyl) N,N,N′,N′-tetraacetic acid (EGTA), N,N′-bis(hydroxybenzyl)ethylenediamine-N,N′-diacetic acid (HBED) and triethylenetrianinehexaacetic acid (TTHA). The chelating groups derived from macrocyclic ligands include, e.g. 1,4,7,10-tetra-azacyclododecane-N,N′,N″,N′″-tetra acetic acid (DOTA), 1,4,7,10-tetra-azacyclotridecane-1,4,7,10-tetra acetic acid (TITRA), 1,4,8,11 -tetra-azacyclotetradecane-N,N′,N″,N′″-tetra acetic acid (TETA), 1,4,8,11 -tetra-azacyclotetradecane (TETRA) and aryl chelating moieties e.g. hydrazinonicotinamide (HYNIC). While conventional chelating agents are within the scope of the present invention, the applicants have, for the first time, employed certain novel MAG3 derivatives as chelating agents. The present invention also encompasses chelating groups based on peptides e.g. preferred derivatives of mercaptoacetyltriglycine (MAG3) which are not previously known to be employed as chelating agents in this field. MAG 3 chelating agents include: SH—CH 2 —CO-Gly-Gly-Gly Cys-Gly-Aib-Ala (SEQ ID NO: 3) Cys-Gly-Gly-Aib (SEQ ID NO:4) Gly-Gly-Ala-Aib (SEQ ID NO: 5) Cys-Aib-Gly-Gly (SEQ ID NO: 6) Cys-Ala-Gly-Aib (SEQ ID NO: 7) Gly-Gly-Gly-Aib (SEQ ID NO: 8) Gly-Gly-Aib-Ala (SEQ ID NO: 9) These MAG3 peptide derivatives are preferredchelating groups. When the MAG3 peptide derivatives are used a spacer group is required. The preferred spacer groups are amino acids of the formula NH 2 —(CH 2 ) n —COOH where n is 4, 5 or 6. When n is 4 the spacer group is 5-amino pentanoic acid. When n is 5 the spacer group is 6-amino hexanoic acid or amino caproic acid. When n is 6 the spacer group is 5-amino heptanoic acid. When a spacer is used, the VIP analog is attached to the carboxylic end of the spacer and the chelating moiety to the amino end. In a preferred embodiment the novel peptide reagent comprises 33 amino acids: 28 from a VIP analog, 1 from a spacer group and 4 from a chelating moiety attached to radiolabeled nuclide to provide a novel and hitherto unknown radiotherapeutic and radioscintographic agent. The methods involved in the synthesis, purification, characterization and radiolabeling of these peptides are illustrated in detail in the examples. The following section also includes biological data relating to the imaging efficacy and dosimetry of the aforesaid radiolabeled peptides. The examples have been furnished for illustrating and providing insight into the invention and should not be construed as limiting the scope of the invention. EXAMPLES Solid Phase Pelptide Synthesis An analog of the present invention can be made by exclusively solid phase techniques, by partial solid phase/solution phase techniques and/or by fragment condensation. Methods for chemical synthesis of polypeptides are well known in the art. Stewart and Young, Solid Phase Peptide Synthesis (W. H. Freeman and Co., 1969), Atherton and Shepherd, 1988, J. Chem. Soc. Perkin Trans. I, 2287. Preferred, semi automated, stepwise solid phase methods for synthesis of peptides of the invention are provided in the examples below. Example 1 Preparation of VIP Analogs Peptides were synthesized using preferably, Fmoc (9-fluorenyl methoxy carbonyl) solid-phase methodology, on CS Bio (Model 536) Peptide Synthesizer (CS Bio Co., San Carlos, Calif., U.S.A.). Sequential assembly of a peptide analog was conducted from the carboxy terminus, by loading of a Fmoc protected amino acid to a solid-phase resin, to the amino terminus. This was proceeded by subsequent removal of the Fmoc protecting group of the amino acid and a stepwise, sequential addition of Fmoc protected amino acids in repetitive cycles to obtain an intermediate protected peptide resin. For peptides that were amidated at the carboxy-terminus, Rink Amide resin was employed, and the loading of the first Fmoc protected amino acid was affected via an amide bond formation with the solid support, mediated by diisopropylcarbodiimide and HOBt. Substitution levels for automated synthesis were preferably between 0.2 and 0.8 mmole amino acid per gram resin. Steps in the synthesis of VIP analogs encompassed in the present invention, employed the following protocol: MIX TIME NO. OF STEP REAGENT (MIN) TIMES 1. Methylene chloride 1 2 2. Dimethyl formamide 1 1 3. 20% Piperidine in 1 1 Dimethyl formamide 4. 20% Piperidine in 29 1 Dimethyl formamide 5. Dimethyl formamide 1 3 6. Isopropanol 1 2 7. Methylene chloride 1 2 8. Amino Acid Variable 1 9. Dimethyl formamide 1 2 10. Stop or Return for next cycle The 9-fluorenyl methoxy carbonyl (Fmoc) group was used for the protection of the α-amino group of all amino acids employed in the syntheses. However, other protecting groups known in the art for α-amino group may be employed successfully. Side chain functional groups were protected as follows: Trityl (trt) and t-butyloxycarbonyl (Boc) were the preferred protecting groups for the imidazole group of Histidine. ydroxyl groups of serine, threonine and tyrosine were protected by t-butyl (t-Bu) groups. mc (2,2,5,7,8-pentamethyl-chroman-6-sulfonyl) and Pbf (2,2,4,6,7-pentamethyldihydro benzofuiran-5-sulfonyl) were the preferred protecting groups for the guanido group in Arginine. Trityl protection was used for asparagine and glutamine. Tryptophan was either used with Boc protection or unprotected. The lysine side chain was Boc protected and aspartic acid and glutamic acid had t-butyl side chain protection. The resin employed for the synthesis of carboxy-amidated analogs was 4-(2′,4′Dimethoxyphenyl-Fmoc-aminomethyl)-phenoxymethyl-derivatized polystyrene 1% divinylbenzene (Rink Amide) resin (100-200 mesh), procured from Calbioichem-Novabiochem Corp., La Jolla, U.S.A., (0.47 milliequivalent NH 2 /g resin). Typically, 2-8 equivalents of Fmoc protected amino acid per resin nitrogen equivalent were used. The activating reagents used for coupling of amino acids in the solid phase synthesis of peptides are well known in the art. These include DCC, DIPCDI, DIEA, BOP, PyBOP, HBTU, TBTU, and HOBt. Preferably, DCC or DIPCDI/HOBt or HBTU/HOBT and DIEA couplings were carried out. Swelling of the resin was typically carried out in dichloromethane measuring to volumes 10-40 ml/g resin. The protected amino acids were either activated in situ or r added in the form of preactivated esters known in the art such as NHS esters, Opfp esters etc. Atherton, E. et al. 1988, J. Chem. Soc. Perkin Trans. I. 2887, Bodansky, M. in “Peptides, Analysis, Synthesis and Biology (E. Gross, J. Meienhofer eds.) Vol. I, Academic Press, New York, 1979, 106. Coupling reaction was carried out in DMF, DCM or NMP or a mixture of these solvents and was monitored by Kaiser test (Kaiser et al., Anal. Biochem., 34, 595-598, 1970). Any incomplete reactions were re-coupled using freshly prepared activated amino acids. After complete assembly of the analog, the amino-terminal Fmoc group was removed using steps 1-6 of the above protocol and then the peptide-resin was washed with methanol and dried. The analogs were then deprotected and cleaved from the resin support by treatment with trifluoroacetic acid, crystalline phenol, ethanedithiol, thioanisole and de-ionized water for 1.5 to 5 hours at room temperature. The crude peptide was obtained by precipitation with cold dry ether, filtered, dissolved and lyophilized. The resulting crude peptide was purified by preperative high performance liquid chromatography (HPLC) using a LiChroCART® C18 (250. Times. 10) reverse phase column (Merck, Darmstadt, Germany) on a Preparative HPLC system (Shimadzu Corporation, Japan) using a gradient of 0.1% TFA in acetonitrile and water. The eluted fractions were reanalyzed on Analytical HPLC system (Shimadzu Corporation, Japan) using a C,8 LiChrospher®, WP-300 (300.Times.4) reverse-phase column. Acetonitrile was evaporated and the fractions were lyophilized to obtain the pure peptide. The identity of each peptide was confirmed by electron spray mass spectroscopy. (a) Preparation of Fmoc-Asn(trt)-Resin A typical preparation of the Fmoc-Asn(trt)-resin was earned out using 0.5 g of 4-(2′,4′-Dimethoxyphenyl-Fmoc-aminomethyl) phenoxymethyl-derivatized polystyrene 1% divinylbenzene (Rink Amide) resin (0.47 mM/g) (100-200 mesh), procured from Calbiochem-Novabiochem Corp., La Jolla, U.S.A. The resin was first allowed to swell in methylene chloride (2. Times. 25ml for 10 min.). It was washed once in dimethylformamide for 1 min. All solvents in the automated protocol were in 20 ml portions per addition. The Fmoc-protecting group on the resin was removed by following steps 3 to 7 of the synthesis protocol. Deprotection of the Fmoc group was checked by the presence of blue beads in a positive Kaiser test. For loading of the first amino acid on the free amino (NH 2 ) group of the resin, the first amino acid, Fmoc-Asn(trt)-OH, was weighed in four fold excess, along with a similar fold excess of HOBt, in the amino acid vessel of the peptide synthesizer. These were dissolved in dimethylformamide (A.C.S. grade) (J. T. Baker, Phillipsburg, N.J., U.S.A.) and activated with DIPCDI, just prior to the addition to the resin in the reaction vessel of the peptide synthesizer. HOBt was added in all coupling reactions, especially in the case of Arg, Asn, Gin and His. The coupling reaction was carried out for a period ranging from 1-3 hours. Loading of the first amino acid was complete when Kaiser test gave a negative result and there was adequate weight increase when the resin, with the first amino acid attached, was dried in vacuum overnight and weighed. Example 2 (b) Synthesis of His-Ser-AsM-Aib-Val-4-Cl-D-Phe-Thr-Asp-Asn-Tyr-Thr-Arg-Leu-Arg-Lys-Gln-Leu-Ala-Val-Lys-Lys-Tyr-Leu-Asn-Ser-Ile-Leu-Asn-NH 2 (SEQ ID NO: 2) The synthesis of SEQ ID NO:2 was initiated by using all of the resin loaded with FmocAsn(trt)-OH as prepared in example (a) above. This was subjected to stepwise deprotection and coupling steps as in steps 1-10 of the synthesis cycle. In each coupling reaction, a four-fold excess of amino-acid, DIPCDI and HOBt were used. The amounts of components are summarized in the table below: CYCLE GRAMS OF PROTECTED AMINO ACID 1. 0.333 Leu 2. 0.333 Ile 3. 0.361 Ser 4. 0.560 Asn 5. 0.333 Leu 6. 0.432 Tyr 7. 0.441 Lys 8. 0.441 Lys 9. 0.319 Val 10. 0.292 Ala 11. 0.333 Leu 12. 0.575 Gln 13. 0.441 Lys 14. 0.625 Arg 15. 0.333 Leu 16. 0.625 Arg 17. 0.374 Thr 18. 0.432 Tyr 19. 0.560 Asn 20. 0.387 Asp 21. 0.374 Thr 22. 0.396 4-Cl-D-Phe 23. 0.319 Val 24. 0.292 Ala 25. 0.387 Asp 26. 0.361 Ser 27. 0.449 His Upon completion of synthesis and removal of the N-terminal Fmoc protecting group (steps 1-6 of the synthesis cycle), the peptide-resin was washed twice with methanol, dried and weighed to obtain 0.560 g. This was subjected to cleavage in a cleavage mixture consisting of trifluoroacetic acid and scavengers, crystalline phenol, ethanedithiol, thioanisole and water for a period of 1.5-5 hours at room temperature with continuous stirring. The peptide was precipitated using cold dry ether to obtain ˜280 mg of crude peptide. The crude peptide was purified on a C 18 preparative reverse phase HPLC column (250. Times. 10) on a gradient of acetonitrile and water in 0.1% TFA as described elsewhere. The prominent peaks were collected and lyophilized, reanalysed on analytical HPLC and subjected to mass analysis. The calculated mass was ˜3356.16 and the mass obtained was 3357.2. The HPLC pure peptide was then subjected to bio-analysis. Incorporation of Spacer/bridging and Chelating Groups to Peptide Derivatives Example 3 Synthesis of DTPA-Spacer-VIP Analog DTPA-Acp-His-Ser-Asp-Aib-Val-4-Cl-D-Phe-Thr-AsD-Asn-Tyr-Thr-Arg-Leu-Arg-Lys-Gln-Leu-Ala-Val-Lys-Lys-Tyl-Leu-Asn-Ser-Ile-Leu-Asn-NH 2 (SEQ ID NO: 10) The attachment of the spacer groups to the peptide derivatives was carried out on solid phase. 0.5 gm of peptide-resin of consisting of peptide sequence (b) (SEQ ID NO:2) was synthesized in the same way as described in Example 2. The N-teminal end was deprotected using piperidine. The spacer Amino caproic acid (Acp) was converted to Fmoc-Acp following the standard method of N-terminal protection of amino acid. Fmoc-Acp (185mg) was dissolved in DMF and coupled to the peptide resin using DIPCDI/HOBt as the coupling agent. The completion of the reaction was monitored by standard Kaiser test. It was further deprotected and coupled to DTPA anhydride in presence of DIPCDI/HOBt. After completion of the reaction it was dried. 530 mg of peptide containing spacer and chelate on resin was obtained and was cleaved as described in Example 2. 282 mg crude peptide-conjugate was obtained which was further purified by HPLC and characterized. Example 4 Synthesis of MAG3 Derivative-Spacer-VIP Analog Cys-Gly-Aib-Ala-Acp-His-Ser-Asp-Aib-Val-4-Cl-D-Phe-Thr-Asp-Asn-Tyr-Thr-Arg-Leu-Arg-Lys-Gln-Leu-Ala-Val-Lys-Lys-Tyr-Leu-Asn-Ser-Ile-Leu-Asn-NH 2 (SEQ ID NO: 11) The spacer group Acp was attached to VIP analog in the same way on solid phase as described in Example 3. 1.3 gm of peptide resin on which VIP analog and spacer group i.e. Acp are assembled was taken and the four amino acids Ala, Aib, Gly and Cys were coupled respectively following the same protocol as described in Example 2. 1.398 of peptide-resin were obtained on cleavage which yielded 675.0 mg of crude peptide conjugate. It was further purified and characterized. General Methods for Radiolabeling In forming a complex of radioactive technetium with a peptide of this invention (SEQ ID NO:2) the technetium complex, preferably a salt of Tc-99m pertechnetate, was reacted with a peptide of this invention in the presence of a reducing agent such as stannous ion, dithionite ion or ferrous ion. In a preferred embodiment, the reducing agent is stannous ion. Example 5 1 mg of the peptide of SEQ ID NO:2 was dissolved in 1 ml of water or 0.9% normal saline. To 100 μg of freshly dissolved peptide 8-15 μg of stannous chloride dissolved in 10% acetic acid was added. pH is set to 5.5 with 0.5N NaHCO 3 , 1-10 mCi of freshly eluted Tc-99m sodium pertechnetate is added to the peptide, the reaction proceeds at room temperature for 15-45 minutes and then is filtered through a 0.22 μm filter. The radiolabeled peptide was either used directly or purified on a Sep Pak C18 cartridge using 50% MeCN-water/0.1% TFA as eluant. The extent of Tc-99m peptide labeling achieved was determined by instant thin layer chromatography (ITLC). 5 μl of the radiopharmaceutical was spotted at the base of silica gel coated ITLC strips and chromatographed with acetone or normal saline. Under these conditions 99% of Tc99m associated radioactivity remained at the origin (Rf=0) in either solvent indicating that no significant concentration of free Tc-99m pertechnetate could be detected in the sample. Example 6 Alternatively, the peptide of SEQ ID NO:2 was reacted with technetium-99m in a reduced form. Example 7 In another alternative, both SEQ ID NO:2 and technetium-99m were reacted with a reducing agent prior to being reacted with each other; preferred reducing agent being stannous ion (other reducing agents include dithionite and ferrous ions). Example 8 In forming a complex of radioactive technetium with the MAG3 chelated peptide (SEQ ID NO:1 1), the technetium complex, preferably a salt of Tc-99m pertechnetate, was reacted with a peptide of this invention in the presence of a reducing agent; in a preferred embodiment, the reducing agent is stannous ion (other reducing agents include dithionite and ferrous ions). 1 mg of the MAG3-peptide was dissolved in 1 ml of water or 0.9% normal saline. To 100 υg of freshly dissolved peptide 8-15 υg of stannous chloride dissolved in 10% acetic acid was added. pH was set to 5.5 with 0.5N NaHCO 3 . 1-10 mCi of freshly eluted Tc-99m sodium pertechnetate was added to the peptide, the reaction proceeded at room temperature for 15-45 minutes and then filtered through a 0.22 μm filter. The radiolabeled peptide was either used directly or purified on a Sep Pak C18 cartridge using 50%MeCN-water/0.1% TFA as eluant. The extent of Tc-99m peptide labeling achieved was determined by instant thin layer chromatography (ITLC). 5 μl of the radiopharmaceutical was spotted at the base of silica gel coated ITLC strips and chromatographed with acetone or normal saline. Under these conditions 99% of Tc99m associated radioactivity remained at the origin (Rf=0) in either solvent indicating that no significant concentration of free Tc-99m pertechnetate could be detected in the sample. Example 9 Alternatively, the MAG3-peptide complex was reacted with technetium-99m in a reduced form. Example 10 In another alternative, both the MAG3-peptide complex and technetium-99m were reacted with a reducing agent prior to being reacted with each other; preferred reducing agent being stannous ion (other reducing agents include dithionite and ferrous ions). Example 11 In forming a complex of radioactive technetium with the DTPA chelated peptide (SEQ ID NO:10), the technetium complex, preferably a salt of Tc-99m pertechnetate, was reacted with the peptide of this invention in the presence of a reducing agent; in a preferred embodiment, the reducing agent being stannous ion (other reducing agents include dithionite and ferrous ions). 1 mg of the DTPA-peptide was dissolved in 1 ml of water or 0.9% normal saline. 100 υg of freshly dissolved peptide added 8-15 μg of stannous chloride dissolved in 10% acetic acid was added. pH was set to 5.5 with 0.5N NaHCO 3 . 1-10 mCi of freshly eluted Tc-99m sodium pertechnetate was added to the peptide, the reaction proceeded at room temperature for 15-45 minutes and then filtered through a 0.22 μm filter. The radiolabeled peptide was either used directly or purified on a Sep Pak C18 cartridge using 50%MeCN-water/0.1% TFA as eluant. The extent of Tc-99m peptide labeling achieved was determined by instant thin layer chromatography (ITLC). 5 μl of the radiopharmaceutical was spotted at the base of silica gel coated ITLC strips and chromatographed with acetone or normal saline. Under these conditions 99% of Tc-99m associated radioactivity remained at the origin (Rf=0) in either solvent indicating that no significant concentration of free Tc-99m pertechnetate could be detected in the sample. Example 12 Alternatively, the DTPA-peptide complex of the invention was reacted with technetium-99m in a reduced form. Example 13 In another alternative, both the DTPA-peptide complex of the invention and technetium-99m were reacted with a reducing agent prior to being reacted with each other; preferred reducing agent being stannous ion (other reducing agents include dithionite and ferrous ions). Example 14 Other radionuclides that may be used to radiolabel the peptides—the VIP receptor antagonists are those known in the art and include 123 I, 131 I, 111 In, and 188 Re etc. In vitro Biological Assays Example 15 Peptides of the invention were assayed for biological activity in homogeneous competition binding assays using 125 I labeled VIP peptide (SEQ ID NO:1) and in heterogeneous displacement assays using either 125 I labeled VIP (10-28) fragment (Tyr-Thr-Arg-Leu-Arg-Lys-Gln-Met-Ala-Val-Lys-Lys-Tyr-Leu-Asn-Ser-Ile-Leu-Asn-NH 2 ) (SEQ ID NO: 12) or 125 I labeled VIP (SEQ ID NO:1). The assays were performed on a variety of human tumor cell lines. In the practice of these methods, the peptide was radioiodinated using the iodogen method. Briefly, 5 μg of the peptide in 10 μl of 50 mM PBS (pH 7.5), an appropriate amount of the radioisotope and 50 μg-100 μg iodogen were incubated at room temperature for 15-30 min with occasional mixing. Radioiodinated peptide was purified from unincorporated radioactive iodine by purification on a Sep Pak C 18 cartridge, essentially following the same procedure as described for technetium labeling. Example 16 Receptor binding and competition assays were performed at 4-8° C. Briefly, 50,000 cells were plated per well of a 24 well plate and allowed to adhere overnight. Before the assay, the cells were washed twice with ice cold binding buffer (25 mM HEPES, 10 mM MgCl 2 and 1% BSA in RPMI 1640 medium). The cells were incubated for 2-3 hrs with an appropriate concentration (0.1-10 nM) of the 125 I labeled peptide (SEQ ID NO:2), in the presence and absence of the cold ligand, which is the uniodinated form of the same peptide (1 nM-10 μM). After incubation, the cells were washed thrice with the binding buffer to remove the unbound peptide. The cells were lysed and counts were measured in a Gamma counter. From a comparison of the extent of binding in the presence or absence of the unlabeled peptide (SEQ ID NO:2), the dissociation constant (Kd) (TABLE I) and maximal binding (Bmax) (TABLE II) were calculated for the peptide. It was observed that the peptide bound to two kinds of receptors on the cell surface. One receptor had a high affinity (nM range) but low surface expression on the cells whereas the other receptor had a low affinity (υM range) but high expression on the cell surface. These characteristics are similar to what has been previously reported for VIP receptors. The following tumor cell lines were assayed using the above described binding competition assay: HT29 (human colorectal adenocarcinoma); PTC (human primary tumor cells adenocarcinoma); KB (human squamous cell carcinoma); 4451 (human squamous cell carcinoma); L132 (human lung carcinoma); A549 (human lung carcinoma); HBL100 (human breast carcinoma) and PA1 (human ovarian carcinoma). Cells were grown in RPMI 1640 supplemented with 10% fetal calf serum, glutamine and antibiotics using standard cell culture techniques (see Animal Cell Culture: A Practical Approach, Freshney, ed, IRL press: Oxford, UK, 1992). TABLE I S No. CELL LINES K d1 (nM) K d2 (μM) 1) PTC 1.77 1.18 2) HT29 3.8 1.37 3) KB 4.2 1.84 4) 4451 4.48 1.91 5) L132 2.6 1.03 6) A549 2.09 1.53 7) HBL 100 2.13 1.6 8) PA1 5.6 1.89 TABLE II SNo. CELL LINES Bmax 1 (M) Bmax 2 (M) 1) PTC 8.3E-10 9.66E-08 2) HT29 3.67E-10 6.72E-08 3) KB 6.02E-10 4.29E-08 4) 4451 5.79E-10 5.1E-08 5) L132 3.06E-10 4.58E-08 6) A549 5.40E-10 4.83E-08 7) HBL 100 6.97E-10 6.91E-08 8) PAI 5.29E-10 4.77E-08 Example 17 Displacement assays were performed at 4-8° C. Briefly, 50,000 cells were plated per well of a 24 well plate and allowed to adhere overnight. Before the assay, the cells were washed twice with ice cold Binding buffer (25 mM HEPES, 10 mM MgCl 2 and 1% BSA in RPMI 1640 medium). The cells were incubated for 2-3hrs with an appropriate concentration ( 0.1-10 nM) of 125 I labeled VIP (10-28) fragment (SEQ ID NO:12) in the presence and absence of the cold ligand (SEQ ID NO: 2) (1 nM-10 μM). The non specific binding was ascertained using 10 μM of VIP. After incubation, the cells were washed thrice with the binding buffer to remove the unbound peptide. The cells were lysed and counts were measured in a Gamma counter. From a comparison of the extent of binding in the presence or absence of the unlabeled peptide, a concentration was determined corresponding to inhibition of 125 I labeled VIP (10-28) fragment binding by 50% (termed the IC 50 ) (TABLE III). The following tumor cell lines were assayed using the above described displacement assay: HT29 (human colorectal adenocarcinoma); PTC (human primary tumor cells adenocarcinoma); KB (human squamous cell carcinoma); 4451 (human squamous cell carcinoma); L132 (human lung carcinoma); A549 (human lung carcinoma); HBL100 (human breast carcinoma) and PA1 (human ovarian carcinoma). Cells were grown in RPMI 1640 supplemented with 10% fetal calf serum, glutamine and antibiotics using standard cell culture techniques (see Animal Cell Culture: A Practical Approach, Freshney, ed, IRL press: Oxford, UK, 1992). TABLE III S No. CELL LINES IC 50 (pM) 1) PTC 132 2) HT29 220 3) KB 350 4) 4451 383 5) L132 228 6) A549 275 7) HBL 100 236 8) PA1 310 Example 18 Displacement assays were performed at 4-8° C. Briefly, 50,000 cells were plated per well of a 24 well plate and allowed to adhere overnight. Before the assay, the cells were washed twice with ice cold Binding buffer (25 mM HEPES, 10 mM 15 MgCl 2 and 1% BSA in RPMI 1640 medium). The cells were incubated for 2-3hrs with an appropriate concentration (0.1-10 nM) of 125 I labeled VIP (SEQ ID NO:1) in the presence and absence of the cold ligand (SEQ ID NO: 2) (1 nM-10 μM). The non specific binding was ascertained using 10 μM of VIP. After incubation, the cells were washed thrice with the Binding buffer to remove the unbound peptide. The cells were lysed and counts were measured in a Gamma counter. From a comparison of the extent of binding in the presence or absence of the unlabeled peptide, a concentration was determined corresponding to inhibition of 125 I labeled VIP binding by 50%. Different tumor cell lines were assayed using the above described displacement assay: HT29 (human colorectal adenocarcinoma); PTC (human primary tumor cells adenocarcinoma); KB (human squamous cell carcinoma); 4451 (human squamous cell carcinoma); L132 (human lung carcinoma); A549 (human lung carcinoma); HBL100 (human breast carcinoma) and PA1 (human ovarian carcinoma). Cells were grown in RPMI 1640 supplemented with 10% fetal calf serum, glutamine and antibiotics using standard cell culture techniques (see Animal Cell Culture: A Practical Approach, Freshney, ed, IRL press: Oxford, UK, 1992). The ligand i.e. VIP analog was able to significantly displace the binding of radiolabelled VIP to the cell lines as shown in Table IV. TABLE IV S No. CELL LINES IC 50 (nM) 1) PTC 2.53 2) HT29 2.78 3) KB 6.2 4) 4451 7.1 5) L132 6.8 6) A549 4.3 7) HBL 100 5.41 8) PA1 7.6 These results demonstrate that the VIP analog described in the invention is capable of specifically binding to VIP receptors in standard in vitro assays on a variety of human tumor cell types. Imaging of Human Tumor Induced in Nude Mice Example 19 Tc-99m labeled peptide (SEQ ID NO: 10) was used to image tumors induced subcutaneously in the abdomen of NIH nu/nu nude mice. Following intravenous administration in human adenocarcinoma tumor bearing nude mice, images were taken at different time intervals post infection, using a conventional gamma camera. A rapid blood clearance was observed with little accumulation in liver and kidney while tumor uptake was found to achieve significant levels as early as 15 min post injection. The major pathway of clearance for the labeled peptide of the invention is through the kidneys as shown by a significant activity in the bladder and urine. These results indicate that the VIP analogue of the present invention (SEQ ID NO:2) has utility as scintigraphic imaging agent for imaging tumor of adenocarcinoma origin in humans. Maximum binding was seen at 3 hours leading to greater accumulation of radioactivity in tumors in comparison to the normal visceral tissue. The results are shown in FIG. 1 which clearly depicts that the accumulation of Tc-99m labeled VIP analog is high in tumor (indicated with arrow) as compared to the accumulation in viscera (unmarked dots) after 3 hours of injection. All publications referenced are incorporated by reference herein including the amino acid sequences listed in each publication. All the compounds disclosed and referred to in the publications mentioned above are incorporated by reference herein, including those compounds disclosed and referred to in the articles cited by the publications. 12 1 28 PRT Sus barbatus 1 His Ser Asp Ala Val Phe Thr Asp Asn Tyr Thr Arg Leu Arg Lys Gln 1 5 10 15 Met Ala Val Lys Lys Tyr Leu Asn Ser Ile Leu Asn 20 25 2 28 PRT Artificial Sequence Description of Artificial Sequence This peptide was synthetically generated 2 His Ser Asp Xaa Val Xaa Thr Asp Asn Tyr Thr Arg Leu Arg Lys Gln 1 5 10 15 Leu Ala Val Lys Lys Tyr Leu Asn Ser Ile Leu Asn 20 25 3 4 PRT Artificial Sequence Description of Artificial Sequence This peptide was synthetically generated 3 Cys Gly Xaa Ala 1 4 4 PRT Artificial Sequence Description of Artificial Sequence This peptide was synthetically generated 4 Cys Gly Gly Xaa 1 5 4 PRT Artificial Sequence Description of Artificial Sequence This peptide was synthetically generated 5 Gly Gly Ala Xaa 1 6 4 PRT Artificial Sequence Description of Artificial Sequence This peptide was synthetically generated 6 Cys Xaa Gly Gly 1 7 4 PRT Artificial Sequence Description of Artificial Sequence This peptide was synthetically generated 7 Cys Ala Gly Xaa 1 8 4 PRT Artificial Sequence Description of Artificial Sequence This peptide was synthetically generated 8 Gly Gly Gly Xaa 1 9 4 PRT Artificial Sequence Description of Artificial Sequence This peptide was synthetically generated 9 Gly Gly Xaa Ala 1 10 29 PRT Artificial Sequence Description of Artificial Sequence This peptide was synthetically generated 10 Xaa His Ser Asp Xaa Val Xaa Thr Asp Asn Tyr Thr Arg Leu Arg Lys 1 5 10 15 Gln Leu Ala Val Lys Lys Tyr Leu Asn Ser Ile Leu Asn 20 25 11 33 PRT Artificial Sequence Description of Artificial Sequence This peptide was synthetically generated 11 Cys Gly Xaa Ala Xaa His Ser Asp Xaa Val Xaa Thr Asp Asn Tyr Thr 1 5 10 15 Arg Leu Arg Lys Gln Leu Ala Val Lys Lys Tyr Leu Asn Ser Ile Leu 20 25 30 Asn 12 19 PRT Sus barbatus 12 Tyr Thr Arg Leu Arg Lys Gln Leu Ala Val Lys Lys Tyr Leu Asn Ser 1 5 10 15 Ile Leu Asn

The present invention encompasses radiolabeled peptide analogs of vasoactive intestinal peptide (VIP) labeled with a radionuclide useful for imaging target sites within mammalian living systems. The invention particularly provides radiolabeled VIP derivatives that bind selectively to the VIP receptor on target cells. Specifically, the invention relates to the radiolabeling of VIP-receptor specific agents and their subsequent use for radiodiagnostic and radiotherapeutic purposes. The invention encompasses methods for radiolabeling these peptides with radio-nuclides and the use of these peptides as scintigraphic imaging agents. The radiolabeled VIP derivatives of the present invention exhibit pharmacological activity and therefore are useful as either imaging agent for visualization of VIP-receptor positive tumors and metastases, as a radiodiagnostic agent or as a radio-therapeutic agent for the treatment of such tumors in vivo by specifically targeting the cytotoxic radionuclide selectively to the tumor site in mammalian living systems.

CROSS REFERENCE TO RELATED APPLICATION [0001] This application is a continuation of Application Ser. No. 10/007,516, filed on Dec. 5, 2001, titled “Filtering Device With Associated Sealing Design And Method”, now U.S. Pat. No. ______, which is incorporated herein by reference in its entirety. FIELD OF THE INVENTION [0002] This invention relates to the field of filtering devices, and more particularly, to a hollow fiber type filter device having a single use or disposable design together with a method for using and manufacturing the same. BRIEF DESCRIPTION OF THE DRAWINGS [0003] FIG. 1 shows an overall side view of an embodiment of the filter device of the present invention. [0004] FIG. 2 shows an exploded isometric view of the filter device of FIG. 1 in an embodiment of the present invention. [0005] FIGS. 3A-3D show various views of the ring of FIGS. 1 and 2 of an embodiment of the filter device of the present invention. [0006] FIG. 4 shows a cross-section view of a portion of the filter device of FIG. 1 in an embodiment of the present invention. DETAILED DESCRIPTION OF THE INVENTION [0007] Referring to the Figures, in which like numerals refer to like portions thereof, FIG. 1 shows an overall side view and FIG. 2 shows an exploded isometric view of an embodiment of the filter device of the present invention. Referring now to FIGS. 1 and 2 , Filter Device 100 in this embodiment of the invention is a dialyzer used for hemodialysis. One skilled in the art will recognize that the filter device of the present invention could also be used for filtering other types of fluids besides blood, including, but not limited to water, sewage, or other types of chemical separation. [0008] Filter Device 100 is a dialyzer utilized by patients with kidney disease who suffer from the adverse effects of toxin build-up in their blood. Dialysis is a process which employs an artificial kidney to remove those toxins. In hemodialysis a dialyzer is used which contains a semipermeable membrane dividing the dialyzer into two compartments. Blood is pumped through one compartment and a dialysate solution through the second. As the blood flows by the dialysis fluid, separated by the semipermeable membrane, blood impurities such as urea and creatinine diffuse through the semipermeable membrane into the dialysis solution by diffusion, convection, and absorption. The electrolyte concentration of the dialysis fluid is set so as to maintain electrolytic balance within the patient. [0009] Dialyzers are known in a variety of configurations. The basic concept is to maximize the surface area of the membrane dividing the blood side from the dialysate side, so that the pressure gradient diffusing toxins from the blood side into the dialysate side and diffusing nutrients or pharmacological agents from the dialysate side into the blood side can operate over a wide area. On the other hand, there are size constraints to the overall three dimensional volume of the device, in order to fit into the hemodialysis apparatus. [0010] Filter Device 100 has a large number of Microfibers 104 (not shown in FIG. 2 ) encased in a Housing 102 , which is a hollow cylinder open at both ends. In other designs, Housing 102 may be open only at one end, and Microfibers 104 are looped in a U-shape in Housing 102 such that both open ends of each microfiber are located at the one open end of Housing 102 (not shown). In either design, thousands of the hollow semipermeable Microfibers 104 carry blood in a pathway through one set of open ends of each Microfiber 104 , through the interior of each Microfiber 104 , and exiting out of the other open end of each Microfiber 104 . [0011] As shown in FIG. 1 , thousands of the hollow semipermeable Microfibers 104 carry blood in a pathway that enters from one end through a first Blood Inlet/Outlet Port 118 to the opposite end and out through a second Blood Inlet/Outlet Port 118 so that blood flows through the interior of each Microfiber 104 in a first direction. Dialysate Inlet/Outlet Ports 110 are also present on opposite ends of Housing 102 . A first Dialysate Inlet/Outlet Port 110 carries dialysate in a pathway into Housing 102 , the dialysate flows through Housing 102 in a countercurrent direction to the blood flow and in the space between each Microfiber 104 , and a second Dialysate Inlet/Outlet Port 110 carries the dialysate out of Housing 102 . The material exchange thus takes place across the semipermeable membrane that is the walls of each Microfiber 104 . Label 114 is preprinted and applied after assembly. A Cap 112 screws into each Blood Inlet/Outlet Port 118 after sterilization, and is utilized to ensure an uncontaminated fluid pathway and is typically not removed until the technician is ready to connect the blood lines. [0012] The design of Filter Device 100 produces a high surface area for material exchange in a relatively low volume device. For example, a Filter Device 100 having a 6.3 cm cylindrical diameter and a 25.4 cm length can easily accommodate a bundle of about 12,000 to 13,000 Microfibers 104 . If each Microfiber 104 has a 0.60 cm circumference and is 24 cm long, the total surface area of all 12,000-13,000 Microfibers 104 is approximately 180 cm 2 . [0013] The manufacture of Filter Device 100 begins by joining Rings 108 into each end of Housing 102 . Each Ring 108 is then joined to Housing 102 . Many different joining techniques may be employed including, but not limited to, spin (friction) welding, laser welding, ultrasonic welding, high frequency welding, gluing, adhesive bonding, solvent bonding, screwing with threads, snap fitting, or any other suitable plastic joining technique. In this embodiment of the invention, spin welding is utilized. A plurality of Nubs 120 spaced apart on the outer surface of Ring 108 constitute the spin welding drive features to assist in the spin welding process. Next, open-ended Housing 102 is filled with a bundle of Microfibers 104 which extend in the longitudinal direction throughout the length of Housing 102 and extending a short distance beyond each end. A Potting Cap 202 ( FIG. 2 ) is attached to each Ring 108 to close off each end of Housing 102 . Housing 102 is then positioned in a centrifuge to allow rotation about an axis perpendicular to the central longitudinal axis, wherein the axis of rotation extends through the midpoint of Housing 102 . Potting Compound 116 is then injected into Dialysate Inlet/Outlet Ports 110 on each end of Housing 102 , is spun in a centrifuge, and the fibers are effectively potted in the dialyzer. Alternatively, each end of Housing 102 may be separately spin welded and injected in a two step process. In one embodiment of the invention, polyurethane is used for Potting Compound 116 . Epoxy or any other suitable compound may also be used as a potting material. The centrifugal force produced by the rotation in the centrifuge forces Potting Compound 116 to each end, where it sets and hardens. [0014] Housing 102 is then removed from the centrifuge, and each Potting Cap 202 is removed from each end to expose the hardened Potting Compound 116 encasing the ends of each Microfiber 104 . Potting Compound 116 and the encased Microfibers 104 at each end are then cut through in a plane perpendicular to the central longitudinal axis of Housing 102 , and the Microfibers 104 longitudinal axes, to expose the interior channels of each Microfiber 104 . The result is that the ends of each Microfiber 104 are open for blood flow through the interior channels of each Microfiber 104 extending through Housing 102 , but the rest of the space surrounding each Microfiber 104 at both ends of Housing 102 is filled with polyurethane, creating a seal between the blood and dialysate. [0015] After the potting and cutting process, a Flange Cap 106 is attached to each Ring 108 and spin welded together, permanently adhering it to Filter Device 100 . This design eliminates an O-ring typically used to assist in the sealing of the blood compartment of a dialyzer. [0016] Dialysate Inlet/Outlet Port 110 in the walls of Housing 102 , which are toward but not at the very ends, remain open for dialysate flow there through. A dialysate line is connected to one Dialysate Inlet/Outlet Port 110 and a dialysate return line is connected to the other Dialysate Inlet/Outlet Port 110 . The dialysate thus flows through the interior of Housing 102 in the space surrounding the Microfibers 104 in one direction. Blood flows from an arterial blood line from a patient connected to a first Blood Inlet/Outlet Port 118 , entering the exposed ends of each Microfiber 104 and flowing through the interior channels through the length of Housing 102 in a countercurrent direction, and then out of the other exposed ends of each Microfiber 104 and back to the patient through a venous blood line connected to a second Blood Inlet/Outlet Port 118 . The blood is thus separated from the dialysate by the semipermeable membranes of the microfiber walls, which allow the transfer of liquids, toxins, and nutrients by solute diffusion and pressure gradients. [0017] Typically, dialyzers are reused. After use in a hemodialysis session for a patient, the dialyzer is cleaned and sterilized for subsequent use by the same patient for a next hemodialysis session. The cleansing, sterilizing, storing, and cataloging of each dialyzer to ensure safe use by the same patient is an expensive and laborious task, and fraught with risk should the dialyzer not effectively have had all of the sterilizing chemicals removed from the dialyzer and the patient be exposed to the sterilizing agent itself. Additionally, if the sterilization process was not able to effectively sterilize the dialyzer, the patient may be subjected to a “non’ biocompatible medical device. Further logistic risk remain in the case the dialyzers get mixed up and the wrong dialyzer is used with the wrong patient. Heretofore, single use dialyzers have been too expensive to manufacture to be very practicable. To accommodate the growing demands of the hemodialysis market for single use or disposable dialyzers, the design of Filter Device 100 of the present invention has solved the high cost problem associated with the current manufacture of disposable dialyzers, but yet maintain the performance and medical requirements necessary for successful hemodialysis. [0018] Various seals in a dialyzer must remain intact, which is of special concern when replacing the currently proven expensive materials, from which many dialyzers are made, with less expensive materials in order to reduce costs. Any dialyzer inherently has at least two sealing regions in its respective design. First, the blood and dialysate compartments must be sealed from each other to ensure that a blood leak does not occur. The second seal consists of sealing either the blood or dialysate compartment from the exterior of the dialyzer. [0019] In nearly all dialyzers currently marketed throughout the world, polyurethane is used as a potting material to seal to the housing to ensure that the blood and dialysate compartments are sealed from each other. An O-ring is typically used to separate the blood from the exterior of the dialyzer. [0020] The seals in a dialyzer must not only maintain their integrity through a specified shelf life duration and during the dialysis treatment process, but must also maintain their integrity during the manufacturing process. [0021] The Filter Device 100 of the present invention utilizes molded parts, including Housing 102 and Flange Caps 106 , made with a polypropylene homopolymer that possess comparable general characteristics to the polycarbonate used in the molded components of the Fresenius Hemoflow series of dialyzers, but is considerably less expensive. The choice of materials for the dialyzer are heavily dependent upon the manufacturing processes employed. Though the optical property of the polypropylene homopolymer is significantly more “hazy” compared to polycarbonate, the blood and dialysate compartments are still readily visible to technicians. [0022] Polyurethane in one embodiment of the invention is used as Potting Compound 116 for Filter Device 100 . Instead of an O-ring, a separate Ring 108 molded from polypropylene is utilized. One Ring 108 is spin welded into each end of Housing 102 . Flange Caps 106 are then spin welded onto Rings 108 after Filter Device 100 has been potted and cut. Other joining techniques as listed above, including laser welding, may be used instead of spin welding. However, spin welding is based on a very simple concept and the process generally can be performed faster, less expensively, and with much less continuous maintenance and re-alignment as compared to laser welding. The weld joint designs utilized in Filter Device 100 are very robust and conducive to the rigors of large scale manufacturing. [0023] During the potting process, the interior portion of each Ring 108 becomes encased in Potting Compound 116 . This creates the first seal between the blood and the dialysate compartments. After potting, the potting caps are removed, the ends are cut, and Flange Caps 106 are spin weld onto Rings 108 . The spin welded region constitutes the second seal region, which seals the blood compartment from the outside of Filter Device 100 (a seal which has typically utilized an O-ring). Filter Device 100 is then conditioned during a low flux conditioning process, and then sterilized. Sterilization may be accomplished in a variety of ways, including ethylene oxide (EtO), steam, or radiation sterilization. [0024] A disadvantage of polypropylene is that its hydrophobic property has a tendency to delaminate from the hydrophilic polyurethane potting material due to the chemistry of surface adhesion between the two materials, resulting in leaks between the blood and dialysate compartments. A two-pronged approach has been taken to solve this delamination problem associated with the use of polypropylene. The first involves building a detailed geometry into the design of Ring 108 to minimize delamination or propagation of the delamination through the creation of physical stops, discussed more fully in relation to FIGS. 3A-3D . The second involves the modification of the surface characteristics of the polypropylene to increase adhesion between it and the polyurethane, also discussed more fully in relation to FIGS. 3A-3D . [0025] FIGS. 3A-3D show various views of an embodiment of the ring of FIGS. 1 and 2 in an embodiment of the single use dialyzer of the present invention. FIG. 3A shows a front view of Ring 108 . FIG. 3B shows a side view of Ring 108 . FIG. 3C shows an isometric cross-sectional view of a portion of Ring 108 as seen along lines B-B of FIG. 3A . FIG. 3D shows a cross-sectional view of Ring 108 as seen along line A-A of FIG. 3A . [0026] Referring now to FIGS. 3A-3D , Ring 108 is shaped to coincide with Housing 102 and Flange Caps 106 that each Ring 108 is mated with. Typically, Housing 102 , Flange Caps 106 , and Rings 108 are circular, but other shapes may also be utilized. Ring 108 has Annular Tongue 316 which fits into an annular groove in Housing 102 formed by Annular Inner Lip 410 and Annular Outer Lip 412 in an interference based snap fit fashion in one embodiment of the invention (see FIG. 4 ). Ring 108 also has Annular Outer Rim 312 and Annular Inner Rim 314 which form an annular groove which is designed to receive Flange Cap 106 in an interference based snap fit (see FIG. 4 ). Potting Cap 202 used in the manufacturing process ( FIG. 2 ) is also designed to fit into this annular groove. [0027] Several methods are available to treat the surface of Ring 108 to modify its surface energy to increase adhesion between it and the polyurethane, including plasma, corona discharge, and flame treatments. By increasing the ability of the surface of Ring 108 to adhere to the polyurethane, Ring 108 has been shown to be effective in eliminating potential issues regarding delamination. A delamination could potentially allow the two fluid pathways to mix outside of the filtering microfibers. The detailed geometry of the design of Ring 108 increases the surface area treatable through surface treatment, enhancing the effects of modifying the surface energy of Ring 108 . [0028] In one embodiment of the invention, a typical surface treatment process which allows for the most practical integration into a clean room automated assembly process is the “corona discharge” surface treatment technique. This treatment method is currently utilized in industry to increase the adhesion of inks, coatings, and adhesives to polyolefins, such as polypropylene. The corona discharge consists of a high voltage electrical discharge that is created between two electrodes across a specified distance. This discharge ionizes the gases present between the electrodes and creates unstable chemical species (mainly free radicals), which possess sufficient energy to initiate bond cleavage at the polymer surface. A small fan is situated just above the corona discharge heads and blows the reactive chemical species onto the polymeric surface of the part being treated, Ring 108 , as shown by arrows 308 in FIG. 3D . Ring 108 is especially well suited to accommodate the corona discharge treatment process, presenting a large surface area due to its geometric design. The corona discharge treatment process is based on the surface being treated to be directly exposed to the electrical discharge, and sections of the surface that are not directly in the “line of sight” of the discharge do not receive as effective treatment. Ring 108 is designed to ensure that the polyurethane interface regions of the ring receive optimal amounts of the surface treatment, while also forcing any delamination that may occur to follow a very difficult pathway. Annular Rounded Ridges 318 on the upper and lower surfaces of Annular Anchor 306 have relatively sharp transitions between them to ensure that optimal amounts of “treatable” area of Ring 108 are exposed to the corona discharge treatment process. When this entire section of Ring 108 is embedded in the Potting Compound 116 , delamination is forced to essentially “start” again and again after being initiated anywhere along the Ring 108 /Potting Compound 116 interface as shown in a close up cross-section of Ring 108 , Housing 102 , and Flange Cap 106 shown in FIG. 4 . The effects of the corona discharge treatment may also be somewhat distributed onto Annular Rounded Ridges 318 in the lower surface of Annular Anchor 306 as the unreacted unstable chemical species will be blown into the center of Ring 108 and react with the lower surface of Ring 108 , which also is embedded in Potting Compound 116 . The thickness of Annular Anchor 306 tends to decrease or taper inwardly from Annular Outer Rim 312 , as opposed to increasing or expanding inwardly, which aids in this surface treatment process. [0029] Covalent bonds are produced on the surface of the polymer as the surface is oxidized during the treatment process. This oxidative coating on the polypropylene surface allows the hydrophilic polyurethane to effectively bond to the modified polypropylene. Because the oxidative coating on the polypropylene has the ability to interact with the oxygen present in the air, and simply the dynamic nature of polymers, the stability of the corona discharge treatment is limited to a specified amount of time. However, once potted, the modified surface of Ring 108 is permanent and does not degrade over time. [0030] A large portion of Ring 108 , Annular Anchor 306 , serves as a mechanical lock and is located at an interior portion of Ring 108 and is completely embedded in Potting Compound 116 . This portion of Ring 108 forces delamination to completely circumvent around and through the Annular Rounded Ridges 318 to create an actual delamination between the blood and dialysate compartments of Filter Device 100 as shown in FIG. 4 . [0031] Another feature of Ring 108 are Radial Channels 302 . As the polyurethane potting mass “backfills” from the ends of Filter Device 100 , the residual air from the ends of Filter Device 100 becomes entrapped due to Annular Rounded Ridges 318 of Annular Anchor 306 portion of the design of Ring 108 . Not allowing the potting mass to bind to the corona discharge treated surface because of an air pocket could potentially create an initiation site for a delamination. To address this situation, Radial Channels 302 are periodically notched perpendicular to Annular Rounded Ridges 318 of the upper surface of Annular Anchor 306 of Ring 108 , which allows the air to escape and not become trapped during “backfilling” of Potting Compound 116 . The upper surface of each Annular Anchor 306 is that surface which faces outward toward the ends of Housing 102 . [0032] FIG. 4 shows a cross-section view of a portion of the single use dialyzer of FIGS. 1 and 2 in an embodiment of the present invention. Referring now to FIG. 4 , Dialysate Compartment 402 and Blood Compartment 404 are the regions of ingress and egress of dialysate and blood through Dialysate Inlet/Outlet Ports 110 and Blood Inlet/Outlet Ports 118 respectively. Annular Inner Lip 410 and Annular Outer Lip 412 of Housing 102 receives Annular Tongue 316 in an interference based snap fit fashion. This connection is spin welded as described above. Typically spin welding of polypropylene does not generally produce extensive spin welding particulate, but material does aggregate around the weld joint in the form of jagged flash (melted polymeric material) which aids in sealing welded parts together. Annular Channel 320 and Annular Channel 408 accommodate the flow of some of the melted flash material that is displaced during the spin welding process. [0033] After the potting and cutting process, in similar fashion Annular Interior Rim 414 and Annular Exterior Rim 416 form an annular groove for receiving Annular Outer Rim 312 of Ring 108 . This connection is spin welded as described above. Annular Channel 406 also accommodates the flow of some of the melted material that is displaced during the spin welding process. Annular Channel 422 is a specially designed area where Flange Cap 106 and Ring 108 seal off against each other during the spin welding process, entrapping additional amounts of melted flash material from the spin welding process. This design insures that no flash material is allowed to invade Blood Compartment 404 . Blood tends to coagulate on any rough surface exposed within Blood Compartment 404 , which would degrade the functioning of Filter Device 100 . One skilled in the art will recognize that Annular Channel 422 will also trap residue material from the other types of joining techniques mentioned above. The flat annular portions seal up against each other and ensure that the flash produced will not be introduced into the blood compartment of Filter Device 100 . However, the welding occurs only at the designated region and not at the flat annular regions where additional amounts of flash may be generated. Additional regions that are designed to contain spin weld flash, or residue material from other types of joining techniques, are located around the Housing 102 /Ring 108 weld interface as Annular Outer Lip 412 extends up from Housing 102 along the exterior of Ring 108 , and around the Flange Cap 106 /Ring 108 weld interface as Annular Exterior Rim 416 extends down from Flange Cap 106 along the exterior of Ring 108 . These areas also minimize the flow of flash, or residue material from other types of joining techniques, outside of Filter Device 100 improving the aesthetic features. [0034] The results of various studies on Filter Device 100 show that the design of Ring 108 provides an excellent surface for the corona discharge treatment prior to potting. Extensive quality and delamination testing from two separate experiments of nearly 600 separate Filter Device 100 samples determined that the current design would have a 0.00% chance of delaminating with an upper binomial confidence level of 0.09%. Extensive testing shows that the design of Filter Device 100 possesses excellent capability of resisting delamination, possesses high performance characteristics, and has significantly reduced manufacturing costs. In addition, the clearance characteristics of Filter Device 100 are among the highest currently available on the market. [0035] Having described the present invention, it will be understood by those skilled in the art that many and widely differing embodiments and applications of the invention will suggest themselves without departing from the scope of the present invention.

A filter device made of less expensive material than comparable filter devices heretofore has basic filter components plus some unique design aspects and an additional ring component. The ring provides an interface inside the filter which enables the potting compound to adhere to the filter and create a seal between a first and second fluid compartment within the filter. An embedded region of the ring possesses a detailed geometry which helps ensure that a delamination would be localized and unable to propagate from the first to the second compartment, maintaining the structural integrity of the filter device. To ensure that the sealing interface remains intact and free from delamination, the ring is subjected to a surface treatment, which modifies the surface energy of the ring. This modified surface energy of the ring allows the hydrophilic potting compound to more effectively bond to the modified hydrophobic ring.

CROSS REFERENCE TO RELATED APPLICATION [0001] This application is a divisional and claims the priority benefit of U.S. patent application Ser. No. 10/691,211 filed Oct. 21, 2003 and entitled “Crosstalk Canceller,” which issued as U.S. Pat. No. 7,263,193 on Aug. 28, 2007; U.S. patent application Ser. No. 10/691,211 is a divisional and claims the priority benefit of U.S. patent application Ser. No. 09/195,745 filed Nov. 18, 1998 and entitled “Crosstalk Canceller,” which issued as U.S. Pat. No. 6,668,061 on Dec. 23, 2003; U.S. patent application Ser. No. 09/195,745 claims the priority benefit of U.S. provisional patent application No. 60/065,637 filed Nov. 18, 1997 and U.S. provisional patent application No. 60/069,015 filed Dec. 10, 1997. The subject matter of the aforementioned applications is incorporated herein by reference. BACKGROUND OF THE INVENTION [0002] This invention pertains to audio signal processing, and specifically to a system and method for crosstalk cancellation. [0003] There are a number of settings in which separate audio signals are prepared for the left and right ears of a listener. Such signals are referred to as binaural signals, and are distinct from stereo signals in that the left and right binaural channels are intended to be heard only by the respective left and right ears of the listener. [0004] Binaural signals are typically used to convey spatial information about the sounds presented. It turns out that a sense of sound source location is created by subtle features imposed on the signals arriving at the left and right ears of the listener [5, 6, 7]. By separately processing left-ear and right-ear signals, as illustrated in FIG. 1 , a sound source can be made to appear at any desired location in a listener's perceptual space. [0005] Such synthetic spatial audio—commonly referred to as 3D audio—has application to video games, teleconferencing, and virtual environments, wherein each sound may be processed so as to appear to originate from its generating object. Another 3D audio application is placing “virtual” speakers about a listener, for instance in a standard home theater surround sound configuration as shown in FIG. 2 . Here, each of five surround signals 30 , 40 , 50 , 60 , 70 is processed according to its location 34 , 44 , 54 , 64 , 74 to form left-ear and right-ear signals 32 , 42 , 52 , 62 , 72 and 33 , 43 , 53 , 63 , 73 , which are used to form the left-ear and right-ear channels 35 and 36 of a binaural signal. Presenting the binaural signal to a listener over headphones gives the impression of a five-speaker surround system, though only the two binaural channels are used. [0006] In all of these applications, headphones or similar transducers are often used to ensure that the left and right binaural channels are delivered, respectively, to the left and right ears of the listener [5, pp. 217-220]. If the binaural signal were played through stereo speakers configured as shown in FIG. 4 , each listener ear would hear both binaural channels. This mixing of the left and right binaural channels, called crosstalk, can significantly degrade the spatial cues in the binaural signal, diminishing the listening experience. [0007] There are, however, situations such as in the case of an arcade game where the use of headphones or earphones is impractical, and it is desired to use stereo speakers to present binaural material. In [1], Atal and Schroeder presented a system called a crosstalk canceler for processing a binaural signal to develop a pair of speaker signals that would deliver the original binaural signal to a properly positioned listener. [0008] The system relies on differences among the transfer functions between the two speakers and the two ears. The basic idea is to cancel the crosstalk appearing in the right ear from the left speaker by sending a negative filtered version of the left speaker signal out the right speaker. The filtering is such that the crosstalk from the left speaker and the canceling signal from the right speaker arrive at the right ear simultaneously as negative replicas of each other, and sum to zero. Left ear crosstalk from the right speaker is similarly eliminated. [0009] The crosstalk canceler proposed in [1] can be very effective, but has several drawbacks which limit its usefulness. First, so that the cancellation signal exactly cancels the crosstalk signal, the listener must be carefully positioned at the so-called sweet spot. In addition, the transition between effective cancellation in the sweet spot and no cancellation out of the sweet spot is very abrupt, making it difficult for listeners to find the sweet spot. Consider a 5 kHz signal having a wavelength of about two inches. The listener only need move his head an inch closer to one speaker than the other to turn the perfect cancellation between the crosstalk and canceling signals into perfect reinforcement between the two. [0010] In addition to restricting listener movement, the canceler [1] is sensitive to the shape of the listener's head and ears. To get effective cancellation, particularly at high frequencies, the canceling signal filter should be tailored to the listener. [0011] The second drawback has to do with the timbre or equalization of the canceled signal as compared to that of the original binaural signal. Listeners in the sweet spot sometimes sense that the canceler output is lacking in low-frequency energy compared to the original binaural signal. Listeners away from the sweet spot complain of phase artifacts and a position sensitive equalization. (Note that the apparent equalization away from the sweet spot is important in some applications. For example, consider a television equipped with stereo speakers and virtual surround sound processing as shown in FIG. 3 . While the crosstalk canceler can deliver the virtual surround binaural signal to listener 80 in the sweet spot, the crosstalk canceler should not compromise the listening experience of those away from the sweet spot.) [0012] To address the restrictions on listener movement, Cooper and Bauck in [2] proposed a crosstalk canceler which cancels only the low frequencies; the high-frequency portion of the binaural input is sent to the output unchanged. Many audio signals have their energy concentrated below a few kilohertz, so that canceling only those frequencies should not significantly diminish the cancellation effect. Because the wavelengths for the canceled portion of the binaural signal are relatively large, the listener has greater freedom of movement before perceiving a change in cancellation effectiveness. Essentially, the canceler trades a less effective cancellation in the sweet spot for a broader sweet spot. [0013] In [3, 4], Cooper and Bauck present a canceler equalization based on the observation that each canceler has a set of so-called “null canceler” frequencies at which the canceling signal filter is orthogonal to—that is, ±90° out of phase from—the direct signal filter. The proposed equalization inverts the sum of the power in the direct and canceling filters at the null canceler frequencies. This equalization is an improvement over the one implied in [1] in that listeners away from the sweet spot hear few artifacts, and those in the sweet spot experience less of a timber change. However, for certain kinds of source material, a timbre change is still noticeable for listeners in and out of the sweet spot. [0014] Therefore it is an object of the present invention to provide a crosstalk canceler allowing greater listener movement while maintaining effective cancellation, and having an equalization which leaves the input binaural signal uncolored. Another object is to develop a canceler which is insensitive to listener head and ear acoustic properties. It is also an object of the present invention to broaden the transition between effective cancellation in the sweet spot and no cancellation outside the sweet spot to help listeners find the sweet spot. Another object of the present invention is to develop a canceler which is relatively free of artifacts away from the sweet spot. Finally, it is an object of the present invention to adapt the equalization to the input signal so as to minimize timbre changes imposed by the canceler. SUMMARY OF THE INVENTION [0015] To provide greater listener freedom of movement, the basic idea is to cancel different frequency bands at different locations, rather than to cancel all frequency bands at the same location as is currently practiced. In this way, changes in listener position do not eliminate cancellation, but shift the part of the signal canceled. In addition, this widening of the sweet spot creates a smooth transition between regions of effective cancellation and no cancellation. [0016] The expectation in canceling different frequency bands at different locations is that while the set of listener positions where some cancellation occurs is broader, the cancellation is everywhere less effective than at the sweet spot of a traditional canceler. That the sweet spot of the new canceler is larger than that of traditional cancelers was verified in listening tests using virtual surround sound, speaker spreader, and one-channel signals as the binaural input. Surprisingly, the inventive canceler was perceived to have nearly as effective cancellation in the sweet spot as the traditional canceler. [0017] In analyzing the signal arriving at a listener's ears from a traditional canceler, it was discovered that unless the listener is precisely positioned, the signal arrives with a timbre change compared to the original binaural signal, irrespective of the cancellation effectiveness. A similar timbre change appears when the acoustic characteristics of the listener's head and ears are not those used in designing the crosstalk canceler, regardless of listener position. [0018] The inventive canceler has an equalization which takes into account the signal arriving at the ears of a variety of listeners positioned in a range of locations. The inventive equalization is the one minimizing the timbre change over an expected range of listener positions and listener acoustic characteristics. Whereas the power spectrum of the traditional crosstalk canceler equalization has a number of peaks and valleys, that of the inventive equalization is by comparison smooth. [0019] The timbre of output from cancelers using the inventive equalization, in fact, is less sensitive to listener position or acoustic properties than is that from the traditional canceler [1]. In addition, the inventive equalization has the unexpected benefit or reducing artifacts for listeners outside the sweet spot. [0020] Finally, it was noted that binaural signals having a large monophonic component seemed to require an equalization with more bass emphasis than did binaural signals with a small monophonic component. Based on this observation, a canceler equalization was developed which depends on the percentage of monophonic signal energy in the input binaural signal. In this way, the canceler equalization may be adapted to the binaural input. [0021] One embodiment of the invention is a crosstalk canceler providing greater listener freedom of movement comprising an input audio signal, two output channels, and a network of filters designed to eliminate crosstalk at the ear of a listener at different listener positions for different frequency bands of the input audio signal. [0022] Another embodiment of the invention is a crosstalk canceler equalization which is less sensitive to listener acoustic characteristics and listener position, said equalization being a spectrally smooth version of an input equalization, the details of which may be optionally determined by anticipated ranges of listener acoustic characteristics and listener positions. [0023] An additional embodiment of the invention is a crosstalk canceler having an equalization designed to leave unchanged at the output the power spectrum of a Gaussian binaural input with a specified crosscoherence. Another aspect of this embodiment is a canceler in which the crosscoherence of the input binaural signal is sensed and used to adapt the characteristics of the canceler. BRIEF DESCRIPTION OF THE DRAWINGS [0024] FIG. 1 shows a synthetic spatial audio display. [0025] FIG. 2 shows a binaural virtual surround sound system. [0026] FIG. 3 shows a stereo speaker virtual surround sound system. [0027] FIG. 4 shows the crosstalk geometry. [0028] FIG. 5 shows a crosstalk canceler. [0029] FIG. 6 shows a lattice crosstalk canceler. [0030] FIG. 7 shows a shuffler crosstalk canceler. [0031] FIG. 8 shows a butterfly crosstalk canceler. [0032] FIGS. 9 a and 9 b show a crosstalk remover example. [0033] FIG. 10 shows an incomplete crosstalk cancellation example. [0034] FIG. 11 shows a crosstalk equalization example. [0035] FIG. 12 shows a crosstalk equalization error example. [0036] FIG. 13 shows an inventive sweet spot position example. [0037] FIG. 14 shows example transfer function ratio magnitudes. [0038] FIG. 15 shows example transfer function ratio phase delays. [0039] FIGS. 16 a and 16 b show an inventive mixing filter example. [0040] FIG. 17 shows sweet spot crosstalk energy. [0041] FIGS. 18 a and 18 b show an inventive mixing filter example. [0042] FIG. 19 shows example sweet spot crosstalk energy. [0043] FIGS. 20 a and 20 b show example inventive residual energy minimizing equalization. [0044] FIG. 21 shows inventive smoothed and interpolated equalizations systems. [0045] FIG. 22 shows a smoothed equalization example. [0046] FIG. 23 shows an interpolated equalization example. [0047] FIG. 24 shows inventive reduced feedback equalization systems. [0048] FIG. 26 shows example inventive equalizations. [0049] FIG. 27 shows a system for adapting crosstalk canceler equalization to signal characteristics. [0050] FIGS. 28 a and 28 b show a system and en example inventive equalization approximation. [0051] FIG. 29 shows a system for mixing filter evaluation. [0052] FIG. 30 shows a system for optimizing sweet spot trajectory. [0053] FIG. 31 shows a system for mixing filter optimization. [0054] FIG. 32 shows a system for computing transfer function means. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENT [0055] For clarity, the invention will be described with respect to the symmetric two-speaker, one-listener crosstalk scenario of FIG. 4 . Modifications needed to apply the invention to asymmetric crosstalk geometries, to multiple listeners, or to more than two speakers will be readily apparent to those skilled in the art. In the following, references to listener position or ear position refer also to listener orientation as well as other geometric factors including speaker position and orientation. In addition, in the following equivalent time-domain and frequency-domain quantities and operations are used interchangeably; any technique discussed or description given in one domain is meant to apply in the other. Finally, the functions “mean” and “average” are to be understood in their general sense, for instance being weighted or unweighted arithmetic, geometric, or trimmed means and the like. [0000] Crosstalk Cancellation [0056] To better appreciate aspects of the present invention, the traditional crosstalk canceler will be described in detail. Referring to FIG. 4 , consider two speakers 100 and 102 symmetrically placed about listener 110 at an angle θ 112 with respect to listener axis 111 . Signals applied to the speakers will arrive at the listener's ears transformed according to near-ear and far-ear transfer functions v(w) 104 and φ(w) 105 embodying, among other effects, the speaker radiation, speaker-listener propagation effects, and acoustic characteristics of the listener. Denoting by s l (t) and s r (t) the left and right speaker signals 101 and 103 , the signals l l (t) 106 and l r (t) 109 appearing at the listener's left and right ears 107 and 108 are given by l l ( t )= v ( t )* s l ( t )+φ( t )* s r ( t ),  (1) l r ( t )=φ( t )* s l ( t )+ v ( t )* s r ( t ),  (2) where * represents convolution, and v(t) and φ(t) are the near-ear and far-ear impulse responses, that is, the inverse Fourier transforms of the near-ear and far-ear transfer functions v(w) and φ(w). Expressed in the frequency domain, the listener ear sound pressure signals are l ( w )= C ( w ) s ( w ),  (3) where l(w) and s(w) are columns containing the listener ear signal and speaker signal Fourier transforms, l ⁡ ( ω ) = [ l l ⁡ ( ω ) l r ⁡ ( ω ) ] , s ⁢   ⁢ ( ω ) = [ s l ⁡ ( ω ) s r ⁡ ( ω ) ] , ( 4 ) and C(w), the crosstalk matrix, contains the speaker-listener transfer functions, C ⁡ ( ω ) = [ v ⁡ ( ω ) ϕ ⁡ ( ω ) ϕ ⁡ ( ω ) v ⁡ ( ω ) ] . ( 5 ) [0057] It is clear that unless the far-ear transfer function φ(w) is zero, a binaural signal applied directly to the speakers will exhibit crosstalk. However, as discussed above, crosstalk may be removed by processing the binaural signal so as to anticipate the changes imposed in propagating from the speakers to the listener. [0058] Consider the processing shown in FIG. 5 . Binaural channels b l (w) 120 and b r (w) 121 are processed by canceler filter network 122 to produce crosstalk canceled speaker signals s l (w) 123 and s r (w) 124 , which, in turn arrive at the ears of the listener transformed by the near-ear and far-ear transfer functions comprising the crosstalk matrix C(w). The listener ear signals l(w) are easily related to the binaural signal b(w), l ( w )= C ( w ) s ( w )= C ( w ) X ( w ) b ( w ),  (6) where b(w) is the column of binaural channel signal transforms, b ⁡ ( ω ) = [ b l ⁡ ( ω ) b r ⁡ ( ω ) ] , ( 7 ) and where the matrix transfer function X(w) is referred to as the canceler matrix. Note that if the inverse of the crosstalk C(w) is realizable, setting the canceler to the crosstalk inverse, X ( w )= C −1 ( w ),  (8) will produce left and right listener ear signals l l (w) 129 and l r (w) 130 equal to the respective input left and right binaural channels b l (w) 120 and b r (w) 121 . [0059] The canceler inverse may be expressed in terms of the near-ear and fare transfer functions, X ⁡ ( ω ) = C - 1 ⁡ ( ω ) = [ v ⁡ ( ω ) - ϕ ⁡ ( ω ) - ϕ ⁡ ( ω ) v ⁡ ( ω ) ] v 2 ⁡ ( ω ) - ϕ 2 ⁡ ( ω ) , ( 9 ) and implemented in the lattice architecture of FIG. 6 . Here, binaural inputs 140 and 141 are applied to filters 142 , 143 , 144 , and 145 , each implementing the transfer function contained in the corresponding element of the canceler matrix (9). The filter outputs are combined to form canceled speaker outputs 152 and 153 . [0060] Note that for the crosstalk inverse to exist, the near-ear and far-ear transfer functions cannot be identical at any frequency. If this were the case, any canceling signal riving at one ear would cancel the original signal in the other ear. Also, note that for X(w) to be realizable, the quantity v 2 (w)−φ 2 (w) needs to be minimum phase. If this is not the case, then its minimum phase equivalent may be used to form its inverse in (9), and the signals appearing in the ear of the listener will be the binaural channel signals shifted in phase by the allpass component of v 2 (w)− 2 (w). [0061] The canceler may also be formed by noting that the crosstalk matrix can be decomposed in terms of the sum and difference of the near-ear and far-ear transfer functions, C ⁡ ( ω ) = 1 2 ⁡ [ 1 1 1 - 1 ] · [ v ⁡ ( ω ) + ϕ ⁡ ( ω ) 0 0 v ⁡ ( ω ) - ϕ ⁡ ( ω ) ] · [ 1 1 1 - 1 ] , ( 10 ) where the diagonalizing matrix F = [ 1 1 1 - 1 ] ( 11 ) is referred to as the shuffler matrix. Noting that the shuffler matrix F is twice its own inverse, the crosstalk canceler X(w) can be written as X ⁡ ( ω ) = ⁢ C - 1 ⁡ ( ω ) = 1 2 ⁡ [ 1 1 1 - 1 ] · [ 1 v ⁡ ( ω ) + ϕ ⁢   ⁢ ( ω ) 0 0 1 v ⁡ ( ω ) - ϕ ⁢   ⁢ ( ω ) ] · ⁢ [ 1 1 1 - 1 ] , ( 12 ) leading to the shuffler canceler architecture shown in FIG. 7 . In this canceler implementation, the sum and difference of binaural input channels 160 and 161 are filtered by shuffler sum filter 164 and shuffler difference filter 165 , respectively, the outputs of which are summed and differenced to form the canceled speaker outputs 170 and 171 . The advantage of this architecture is that only two filters are needed, rather than the four required by the lattice canceler shown in FIG. 6 . [0062] The crosstalk inverse may also be decomposed as follows, C - 1 ⁡ ( ω ) = [ 1 - ρ ⁢   ⁢ ( ω ) - ρ ⁢   ⁢ ( ω ) 1 ] · 1 v ⁡ ( ω ) · 1 1 - ρ 2 ⁡ ( ω ) , ( 13 ) where p(w) is the ratio of the far-ear transfer function to the near-ear transfer function, p ( w )=φ( w )/ v ( w ).  (14) The corresponding canceler may be implemented in two stages using the butterfly architecture shown in FIG. 8 . The first stage 192 is referred to as the crosstalk remover or mixing stage, and adds to each binaural channel a filtered version of the other binaural channel; its transfer function is given by R ⁡ ( ω ) = [ 1 - r ⁢   ⁢ ( ω ) - r ⁢   ⁢ ( ω ) 1 ] , ( 15 ) where r(w) is referred to as the mixing filter. The second stage 193 , which may be applied either before or after the first stage, equalizes the output, and is called the canceler equalization; its transfer function is Q ( w )= q ( w ) I,   (16) where I is the identity matrix, and q(w) is the equalization filter. By setting the mixing filter to the transfer function ratio r ( w )= p ( w ),  (17) and the equalization filter to the product q ( w )=1/[ v ( w )(1− p 2 ( w ))],  (18) the butterfly architecture of FIG. 8 will implement the canceler inverse. [0063] To understand the function of the mixing stage R(w), consider the example shown in FIG. 9 . Binaural signal channels 200 and 201 are applied to mixing stage 202 , which produces speaker signals 207 and 208 in response. These signals propagate to the listener, appearing as listener ear signals 215 and 216 . For purposes of illustration, the near ear transfer function here is one v(w)=1, and the far-ear transfer function is a scaled pure delay φ(w)=pe −jwr . In this example, the mixing filter r(w) is set to the transfer function ratio p(w)=φ(w)/v(w)=pe −jwr . [0064] Referring to FIG. 9 , pulse 230 applied to the left binaural channel appears directly at the left speaker as pulse 232 . It also appears delayed and scaled according to −p(w) at the right speaker as pulse 235 . The listener left ear will hear pulse 232 directly from the left speaker via near-ear transfer function 211 v(w)=1. The left ear will also hear pulse 235 , delayed and scaled according to far-ear transfer function 213 φ(w)=pe −jwr . The listener right ear will hear pulse 232 from the left speaker via far-ear transfer function 212 , and pulse 235 directly via near-ear transfer function 214 . [0065] Note that pulses 241 and 242 arriving at the right ear cancel. Pulse 241 arriving from the left speaker via far-ear transfer function 213 is delayed and scaled by the same amount as pulse 235 by mixing filter 203 and near-ear transfer function 214 . Therefore, signals applied to left binaural input 200 do not appear at the listener's right ear. Similarly, right binaural channel signals will be canceled at the listener's left ear. More generally, when the mixing filter r(w) is set to the ratio of the near-ear and far-ear transfer functions, binaural signals processed according to the mixing stage (15) will appear at the listener's ears without crosstalk. [0066] Note that listener ear signals 215 and 216 are not the original binaural signal channels 200 and 201 ; each ear contains an echo of its respective binaural channel 239 and 243 as a residual effect of canceling crosstalk. The purpose of the equalization is now clear: In addition to inverting the near-ear transfer function (referred to as “naturalization” in [3, 4]), the equalizer must eliminate the echo. As shown in FIG. 11 , the echo at the listener ear may be removed by adding a series of echoes to the binaural signal. If the echoes are properly spaced in time and filtered, then the chain binaural signal echoes arriving from the far speaker will exactly cancel all but the first of the binaural signal instances arriving directly from the near speaker. [0000] Inventive Crosstalk Removal [0067] The canceler sensitivity to listener position and listener acoustic characteristics discussed above is seen to result from discrepancies between the mixing filter r(w) and the transfer function ratio p(w). As illustrated in FIG. 10 , the crosstalk signal is the crosstalk binaural channel (i.e., the left binaural channel at the right ear or the right binaural channel at the left ear) filtered by φ(w)−r(w)v(w). As the listener moves, the transfer functions φ(w) and v(w) change, and, unless those changes are anticipated by the mixing filter r(w), the canceling signal radiated from the near-ear speaker will not cancel crosstalk from the far-ear speaker. [0068] To give the listener some freedom of movement while maintaining effective (though not complete) crosstalk cancellation, Cooper and Bauck set the mixing filter to a low-pass filtered version of the transfer function ratio, r(w)=p(w)h(w), h(w) being a low-pass filter with a cutoff frequency above 600 Hz and below 10 kHz. In doing so, crosstalk is canceled only below the cutoff frequency. However, since low frequencies have relatively long wavelengths, p(w) is somewhat insensitive to listener position at low frequencies. As a result, the listener is afforded a degree of freedom of movement without noticeably changing canceler effectiveness. [0069] The present invention gives the listener freedom of movement by canceling different frequency bands at different listener positions. For instance, low frequencies might be canceled at a speaker separation angle of θ=10°, and high frequencies at an angle of θ=30°. Doing so provides a measure of cancellation over a range of anticipated listener positions; listener position changes do not eliminate cancellation, but simply shift the part of the signal canceled. An additional benefit of distributing the cancellation location is that a smooth transition between regions of effective cancellation and no cancellation is created. [0070] Changing the cancellation geometry as a function of frequency may be accomplished by setting the mixing filter to the transfer function ratio evaluated at a frequency-dependent geometry as shown in FIG. 29 , r ( w )= p ( w ,θ( w )),  (19) where θ(w), called the sweet spot trajectory, specifies the frequency-dependent crosstalk geometry at which the transfer function ratio is evaluated. The mixing filter thus designed can be implemented directly as mixing filter 182 and 183 in mixing stage 192 of the butterfly canceler in FIG. 8 . It can also be used in forming the canceler matrix X(w), and implemented as a lattice, shuffler, or other canceler. Equivalently, shuffler or lattice cancelers, (12) or (9), or other cancelers, may be designed directly based on a frequency-dependent geometry. [0071] Details of the sweet spot trajectory θ(w) depend on, among other factors, the desired listener and speaker positions, and the binaural source material. In one embodiment, shown in FIG. 13 , the sweet spot center is moved further from the speakers with increasing frequency. By changing the sweet spot center location more rapidly with decreasing frequency, this embodiment attempts to maintain a constant, but acceptable, level of crosstalk within the extended sweet spot. In another embodiment, the magnitude and phase of the mixing filter are determined from separate sweet spot center trajectories. [0072] In FIG. 14 and FIG. 15 , example transfer function ratio magnitudes and phase delays are shown as functions of frequency for listener positions along the listener axis. Mixing filters based on the inventive sweet spot trajectory 280 and prior art constant sweet spot trajectories 281 , 282 are shown in FIG. 16 . Note that the inventive mixing filter takes on the characteristics of the closer prior art filter at low frequencies and those of the farther prior art filter at high frequencies. [0073] The total energy in the crosstalk signal at an ear of a listener positioned at θ is given by E c (θ)=∫ 0 x |v ( w ,θ) r ( w )−φ( w ,θ)| 2 dw,   (20) where v(w,θ) and φ(w,θ) are the near-ear and far-ear transfer functions to the ear of the listener at θ. The crosstalk energy is plotted in FIG. 17 for the mixing filters implied by the sweet spot center trajectories of FIG. 13 . Note that the inventive sweet spot 300 is somewhat more extended than that of the prior art canceler 301 (corresponding to constant sweet spot 281 ), and of comparable extent to that of prior art canceler 302 (corresponding to constant sweet spot 282 ). [0074] In another embodiment of the invention, the sweet spot trajectory θ(w) is designed to maximize the area over which the listener can move while maintaining a minimum level of crosstalk rejection or maximum level of uncanceled crosstalk energy. In another embodiment, θ(w) is chosen to minimize the maximum crosstalk energy experienced by a listener located in a given region. In optimizing the sweet spot trajectory θ(w) as shown in FIG. 30 , note that it may be useful to weight the crosstalk energy in frequency or position to give more importance to certain spectral bands or listener positions, or to account for the canceler equalization. For instance, the power spectrum of many sounds approximates a 1/w characteristic away from DC, so that in optimizing the sweet spot trajectory, it is useful to weight the crosstalk energy away from DC by 1/w. [0075] Another approach shown in FIG. 31 is to find the optimal mixing filter directly, rather than using θ(w) to parameterize the solution. In this embodiment of the invention, the crosstalk energy is written in terms of the mixing filter and the near-ear and far-ear transfer functions at each frequency and crosstalk geometry of interest, E c (θ, w )=γ( w )·| v ( w ,θ) r ( w )−φ( w ,θ)| 2 ,  (21) where γ(w) represents the product of the equalization filter power and the anticipated signal power at frequency w. The mixing filter r(w) is then taken to be the one optimizing some aspect of the crosstalk energy E c (θ,w). One choice is to minimize the maximum weighted energy over some set of canceler geometries or listener characteristics, r ^ ⁡ ( ω ) = Arg [ min r ⁡ ( ω ) ⁢ { max θ ∈ Θ ⁢ { ∫ 0 π ⁢ w ⁡ ( θ , ω ) ⁢ E c ⁡ ( θ , ω ) ⁢   ⁢ ⅆ ω } } ] , ( 22 ) where w(θ,w) is a weighting reflecting the importance of eliminating crosstalk energy at frequency w and geometry θ, and θ represents the range of canceler geometries and listener characteristics under consideration. Another choice is to maximize the area over which the weighted crosstalk energy is less than a given level, r ^ ⁡ ( ω ) = Arg [ max r ⁡ ( ω ) ⁢ { ∫ θ ⁢   ∈ Θ ⁢ 1 ⁢ ( ∫ 0 π ⁢ w ⁡ ( θ , ω ) ⁢ E c ⁡ ( θ , ω ) ⁢   ⁢ ⅆ ω < v ⁡ ( θ ) ) ⁢   ⁢ ⅆ θ } ] , ( 23 ) where 1(•) is an indicator function, taking on a value of 1 if the condition is true and 0 otherwise, and the quantity v(θ) specifies the maximum acceptable crosstalk energy level as a function of position. Alternatively, the maximum acceptable crosstalk energy level could depend on frequency as well as position, r ^ ⁡ ( ω ) = Arg [ max r ⁡ ( ω ) ⁢ { ∫ θ ⁢   ∈ Θ ⁢ ∫ 0 π ⁢ 1 ⁢ ( E c ⁡ ( θ , ω ) < v ⁡ ( θ , ω ) ) ⁢   ⁢ ⅆ ω ⁢   ⁢ ⅆ θ } ] . ( 24 ) Still another optimization choice is to find the mixing filter minimizing the total crosstalk energy in a given region, r ^ ⁡ ( ω ) = Arg [ min r ⁡ ( ω ) ⁢ { ∫ θ ⁢   ∈ Θ ⁢ ∫ 0 π ⁢ w ⁡ ( θ , ω ) ⁢ E c ⁡ ( θ , ω ) ⁢   ⁢ ⅆ ω ⁢   ⁢ ⅆ θ } ] , ( 25 ) where the weighting w(θ,w) weights the importance of having effective cancellation at a given frequency and speaker-listener geometry. [0076] As an example, FIG. 18 shows the magnitude 450 and phase delay 460 of the prior art mixing filter designed to cancel crosstalk at the ears of a listener positioned on the listener axis twice as far from the line joining the speakers as the distance separating the speakers. Also shown are the magnitude and phase delay of the filter minimizing the total crosstalk energy (25) 451 , 461 and minimizing the maximum crosstalk energy (22) 452 , 462 for listeners on the listener axis between 1.5 and 2.5 times the speaker separation from the speaker axis. Note that magnitude of the optimal mixing filters is similar to that of prior art mixing filters for listener positions closer to the speakers than that used to generate prior art mixing filter magnitude 450 . By contrast, the phase delay of the inventive mixing filters is more like that of prior art mixing filters associated with positions further from the speakers than that used to form prior art mixing filter phase delay 460 . The crosstalk energy associated with the inventive and prior art mixing filters of FIG. 18 is plotted as a function of position in FIG. 19 . The minimizer of the maximum crosstalk energy over the region 452 , 462 provides the widest sweet spot 472 . The prior art crosstalk has the smallest sweet spot 470 and the most abrupt transition between regions of effective cancellation and little cancellation. [0077] Another optimization choice is suggested by the observation that listeners prefer cancelers having a gentle transition between areas of effective cancellation and no cancellation over cancelers with a more abrupt transition. To accommodate this preference, the mixing filter may be optimized so that the slope (derivative with respect to position) of the crosstalk energy in the transition region is minimized. [0078] It should be noted that the optimal mixing filter {circumflex over (r)}(w) (25) may be expressed in closed from, r ^ ⁡ ( ω ) = μ ϕ ⁡ ( ω ) ⁢ μ v ⁡ ( ω ) * + σ ϕ ⁢   ⁢ v * ⁡ ( ω ) μ v ⁡ ( ω ) ⁢ μ v ⁡ ( ω ) * + σ v ⁢   ⁢ v * ⁡ ( ω ) , ( 26 ) where •* denotes complex conjugation, μ φ (w) and μ v (w) are the near-ear and far-ear transfer function means over position, μ φ ( w )=∫ w (θ, w )φ( w ,θ) dθ,   (27) μ v ( w )=∫ w (θ, w ) v ( w ,θ) dθ,   (28) and σ vv *(w) and σ φv *(w) are variances over position, σ vv *( w )=∫ w (θ, w )| v ( w )−μ v ( w )| 2 dθ,   (29) σ φv *( w )=∫ w (θ, w )[φ( w )−μ φ ( w )][ v ( w )−μ v ( w )]* dθ.   (30) Note that the optimal mixing filter has a magnitude and phase approximating that of the mean over position of the transfer function ratio p(w,θ), with the magnitude reduced at frequencies where the transfer function ratio changes rapidly with position. This motivates another embodiment of the invention shown in FIG. 32 , wherein the magnitude or phase of the mixing filter is given by the respective means over position of the magnitude or phase of the transfer function ratio filter, possibly reducing the mixing filter magnitude at any selected frequency by an amount dependent on the transfer function ratio position variance (i.e., the sensitivity of the transfer function ratio to changes in listener position) at that frequency. Inventive Equalization [0079] Listener freedom of movement is also restricted by the canceler equalization. As illustrated in FIG. 11 , the equalization associated with the crosstalk matrix inverse removes the unwanted binaural signal echo by creating two chains of canceling echoes. Unfortunately, as shown in FIG. 12 , the resulting listener ear signals are very sensitive to listener position, which determines the relative alignment and strength of the two chains through the near-ear and far-ear transfer functions. [0080] What is needed is to balance the desire to maintain the original binaural signal equalization with the need to accommodate varying crosstalk geometries and listener characteristics. The inventive canceler equalization achieves this balance by optimizing the equalization over a set of anticipated listener positions and characteristics. This approach differs from that of the prior art which uses a single crosstalk geometry in designing the canceler equalization. [0081] The binaural channel signal appearing at the ear of the listener is filtered by q(w)(v(w,θ)−φ(w,θ)r(w)), q(w) being the canceler equalization filter, r(w) the canceler mixing filter, and v(w,θ) and φ(w,θ) the near-ear and far-ear transfer functions evaluated at the crosstalk geometry and listener characteristics θ. Ideally, the binaural channel would appear at the listener unfiltered; the energy in the difference between the unit transfer function and that imposed on the binaural channel, called the equalization residual is given by E q ( w ,θ)=| q ( w )( v ( w ,θ)−φ( w ,θ) r ( w ))−1| 2 .  (31) [0082] In one embodiment of the invention, the equalization q(w) is optimized to minimize the equalization residual E q (w,θ) over a distribution of crosstalk geometries and listener characteristics p(θ), q ^ ⁡ ( ω ) = Arg [ min q ⁡ ( ω ) ⁢ { ∫ θ ⁢   ∈ Θ ⁢ ∫ 0 π ⁢ ρ ⁡ ( θ ) ⁢ E q ⁡ ( θ , ω ) ⁢   ⁢ ⅆ ω ⁢   ⁢ ⅆ θ } ] , ( 32 ) This solution is available in closed form, q ^ ⁡ ( ω ) = ∫ ρ ⁡ ( θ ) ⁢ ( v ⁡ ( ω , θ ) - ϕ ⁡ ( ω , θ ) ⁢ r ⁡ ( ω ) ) ⁢   ⁢ ⅆ θ ∫ ρ ⁡ ( θ ) ⁢  v ⁡ ( ω , θ ) - ϕ ⁡ ( ω , θ ) ⁢ r ⁡ ( ω )  2 ⁢   ⁢ ⅆ θ . ( 33 ) Denoting by μ v (w) and μ φ (w) the means of the near-ear and far-ear transfer functions with respect to p(θ), μ φ ( w )=∫ p (θ)φ( w ,θ) dθ,   (34) μ v ( w )=∫ p (θ) v ( w ,θ) dθ,   (35) and by σ vv *(w), σ φφ *(w), and σ φv *(w) the variances with respect to p(θ) σ vv *( w )=∫ p (θ)| v ( w )−μ v ( w )| 2 dθ,   (36) σ φφ *( w )=∫ p (θ)|φ( w )−μ φ ( w )| 2 dθ,   (37) σ φv *( w )=∫ p (θ)[φ( w )−μ φ ( w )][ v ( w )−μ v ( w )]* dθ,   (38) the optimal equalization may be written as q ^ ⁡ ( ω ) = 1 μ v ⁡ ( ω ) · 1 1 - r ⁡ ( ω ) ⁢ μ ϕ ⁡ ( ω ) / μ v ⁡ ( ω ) + [ σ vv * ⁡ ( ω ) +  r ⁡ ( ω )  2 ⁢ σ ϕ ⁢   ⁢ ϕ * ⁡ ( ω ) - 2 ⁢ ℜ ⁢   ⁢ { r ⁡ ( ω ) ⁢ σ ϕv * ⁡ ( ω ) } μ v ⁡ ( ω ) ⁢ μ v ⁡ ( ω ) * ⁢ ( 1 - r ⁡ ( ω ) ⁢ μ ϕ ⁡ ( ω ) / μ v ⁡ ( ω ) ) ] , ( 39 ) where {•} is the real part of its argument. By comparison to the prior art equalization, q ⁡ ( ω ) = 1 v ⁡ ( ω ) · 1 1 - r ⁡ ( ω ) ⁢ ϕ ⁢   ⁢ ( ω ) / v ⁡ ( ω ) , ( 40 ) the optimal equalization (39) generates similar train of echoes, but with a shorter time constant (since the bracketed term is nonnegative), particularly in those parts of the spectrum where the near-ear and far-ear transfer functions are sensitive to position changes. In the frequency domain, the magnitude of the optimal equalization will appear smoothed relative to that of the prior art equalization. Note that the greater the sensitivity to position changes or listener characteristics exhibited by v(w) and φ(w), or the greater the range of expected geometries and listeners p(θ), the more smoothed the optimal equalization magnitude compared to the prior art equalization. [0083] As an example, FIG. 20 shows the prior art equalization magnitude 340 along with that of two optimal equalizations. Equalization 341 is designed to minimize the expected equalization residual for listeners uniformly distributed on the listener axis between 1.5 and 2.5 times the speaker separation distance from the speaker axis; equalization 342 minimizes the equalization residual for listeners between 1.0 and 2.5 times the speaker separation from the speaker axis. The equalization residual as a function of listener position is also shown in FIG. 20 . The inventive equalization residuals 344 , 345 achieve their minima over wider ranges of listener position than does the prior art equalization residual 343 . In addition, away from the sweet spot center, the inventive equalization residuals are smaller than the prior art equalization residual. [0084] The observation that the optimal equalization magnitude is essentially a smoothed version of the prior art equalization magnitude leads to the inventive equalizations shown in FIG. 21 and FIG. 24 . In the embodiment shown in FIG. 21 , the inventive canceler equalization spectrum is a smoothed or interpolated version of the spectrum of an input canceler equalization. Note that the smoothing or interpolation may be applied to the entire spectrum, or may be restricted to all but the naturalization, 1/|v(w)| 2 . A smoothed canceler equalization spectrum may be found by applying a running mean (arithmetic, geometric, trimmed or other means may be applied) to a prior art equalization spectrum  q ⁡ ( ω )  2 = 1  v ⁡ ( ω )  2 · 1 1 +  r ⁡ ( ω ) ⁢ ϕ ⁢   ⁢ ( ω ) / v ⁡ ( ω )  2 - 2 ⁢   ⁢ ℜ ⁢ { r ⁡ ( ω ) ⁢ ϕ ⁡ ( ω ) / v ⁡ ( ω ) } . ( 41 ) It may be equivalently found as the spectrum associated with the appropriately windowed version of the prior art equalization impulse response. In FIG. 22 , example prior art equalization 350 is shown along with inventive smoothed equalizations 351 , 352 . Smoothed equalizations 351 , 352 were formed by critical band smoothing of the prior art power spectrum using smoothing bandwidths of 1.0 and 2.0 critical bands, respectively. [0085] An interpolated spectrum may be found by interpolating in the prior art equalization power spectrum points where the quantity r(w)φ(w)/v(w) achieves the same phase. The resulting power spectrum is given by  q ^ ⁡ ( ω )  2 = 1  v ⁡ ( ω )  2 · 1 1 +  r ⁡ ( ω ) ⁢ ϕ ⁢   ⁢ ( ω ) / v ⁡ ( ω )  2 - 2 ⁢   ⁢ α ⁢  r ⁡ ( ω ) ⁢ ϕ ⁡ ( ω ) / v ⁡ ( ω )  , ( 42 ) where αε[−1,1] which determines the points of the prior art equalization interpolated. Several example interpolated equalization magnitudes 361 , 362 are plotted in FIG. 23 along with the prior art equalization magnitude 360 ; interpolation points 363 are marked. [0086] The embodiment of FIG. 24 augments a prior art canceler equalization implementation with an additional filter α(w) which has the effect of reducing feedback, thereby smoothing the spectrum of the prior art canceler. So as to approximate the optimal equalization, feedback should be preferentially reduced in those frequency bands where the feedback is largest. In one instance, a filtered version of the output is added to the feedback path of the prior art equalization, q ^ ⁡ ( ω ) = 1 v ⁡ ( ω ) · 1 1 - r ⁡ ( ω ) ⁢ ϕ ⁡ ( ω ) / v ⁡ ( ω ) + α ⁡ ( ω ) , ( 43 ) where α(w) is a filter having a phase generally similar to that of r(w)φ(w)/v(w); it's presence selectively reduces decay time. In another instance, feedback is reduced directly, q ^ ⁡ ( ω ) = 1 v ⁡ ( ω ) · 1 1 - α ⁡ ( ω ) ⁢ r ⁡ ( ω ) ⁢ ϕ ⁡ ( ω ) / v ⁡ ( ω ) , ( 44 ) where α(S) is a filter (preferably minimum phase) having a magnitude no greater than one; it reduces decay time by limiting the amount of feedback at any given frequency. Note that it is possible to adjust both instances of α(w) above so that the resulting equalization approximates the optimal equalization (39). [0087] Another consideration in crosstalk canceler equalization is the apparent coloring of the binaural signal experienced by those listeners outside the sweet spot. To minimize equalization artifacts for these listeners, the approach taken here is to equalize the canceler so as to be compatible with—i.e., pass unchanged in equalization—certain classes of input signals. For example, many signals including virtual surround binaural signals have a large fraction of their energy common to both binaural channels. In this case, a crosstalk canceler equalized to pass unchanged monophonic signals would be appropriate. The response of a crosstalk canceler X(w)=q(w)R(w) to a two-channel monophonic signal b(w)=m(w)1 is s ( w )= q ( w )(1− r ( w )) m ( w )1.  (45) Setting the equalization to q ⁡ ( ω ) = 1 1 - r ⁡ ( ω ) ( 46 ) leaves the canceler output equal to the canceler input for monophonic inputs. [0088] Consider a binaural input b(w) composed of zero-mean Gaussian random processes having identical power spectra P b (w) and crosscoherence η, E ⁢ { b ⁡ ( ω ) ⁢ b ⁡ ( ω ) T } = P b ⁡ ( ω ) ⁡ [ 1 η η * 1 ] , ( 47 ) where E{•} is the expectation operator and • T is the Hermetian transpose. (Note that the binaural channel crosscoherence η is the energy in the product of the binaural channel signals normalized by the mean of the Individual channel signal energies, so that it takes on values in the range [−1,1]. The energies, and therefore η, may be evaluated as functions of frequency, or they may represent the total energy over the band.) The total power appearing at the output of a canceler X(w)=q(w)R(w)—the sum of the left and right channel output powers—in response to the Gaussian input b(w) is E{s ( w ) T s ( w )}=2| q ( w )| 2 P b ( w )(1+| r ( w )| 2 −2 {ηr ( w )}).  (48) Accordingly, the inventive equalization has a power given by  q ⁡ ( ω )  2 = 1 1 +  r ⁡ ( ω )  2 - 2 ⁢ ℜ ⁢ { η ⁢   ⁢ r ⁢   ⁢ ( ω ) } , ( 49 ) so as to leave the total power of a random process with channel crosscoherence η unchanged at the output. It is worth pointing out that if the input binaural signal were a deterministic signal decomposed into sum—that is, monophonic—and difference components, with η measuring the percentage monophonic energy less the percentage difference energy, the equalization (49) leaves the total output power unchanged. [0089] Note that if the input were monophonic, the channel crosscoherence η would be one, and the equalization power would be that of the monophonic compatible equalization above,  q ⁡ ( ω )  2 ⁢ 1 1 +  r ⁡ ( ω )  2 - 2 ⁢ ℜ ⁢ { r ⁡ ( ω ) } . ( 50 ) If the input channels were statistically independent, the channel crosscoherence would be zero, and the inventive equalization power would be  q ⁡ ( ω )  2 = 1 1 +  r ⁡ ( ω )  2 . ( 51 ) The inventive equalization magnitude is plotted in FIG. 26 for a range of binaural channel crosscoherence values η. [0090] In many cases, the channel crosscoherence will be approximately known a priori. For instance, movie soundtracks presented in binaural virtual surround sound format as shown in FIG. 3 typically have a channel crosscoherence in the range ηε[0.8,0.9]. In one embodiment, if the channel crosscoherence is not known a priori, the listener may tune the canceler equalization to his liking by adjusting the channel crosscoherence value used to determine the equalization power. In another embodiment, shown in FIG. 27 , the binaural channel crosscoherence is sensed (possibly as a function of frequency) and used to adjust the canceler equalization. Alternatively, the percentage of sum and difference energies may be used to set η. [0091] Because of the manner in which the equalization power (49) depends on the binaural channel crosscoherence η, it is difficult to adapt the equalization filter to real-time changes in η. However, the embodiment of FIG. 28 shows an equalization filter comprising two filters in a feedback delay network which has a magnitude approximating that of (49). By setting the delay to the near-ear-far-ear arrival time difference implied by the mixing filter r(w), and by designing the filters α(w) and β(w) to have magnitudes that approximate  α ⁡ ( ω )  = γ - [ γ 2 - 1 ] 1 2 , γ = 1 +  r ⁡ ( ω )  2 2 ⁢  η ⁢   ⁢ r ⁡ ( ω )  ( 52 )  β ⁡ ( ω )  = [ 1 +  α ⁡ ( ω )  2 1 +  r ⁡ ( ω )  2 ] 1 2 , ( 53 )

The invention is a crosstalk canceler wherein different frequency bands are canceled at different locations so as to allow greater listener movement about the “sweet spot” while maintaining effective crosstalk cancellation. A spectrally smooth canceler equalization is used, reducing artifacts for listeners away from the sweet spot and further enlarging the sweet spot. Finally, the canceler equalization is adapted to either the anticipated or the actual crosscoherence among the input channels, producing a natural equalization regardless of the input.

STATEMENT REGARDING FEDERALLY SPONSORED RESEARCH OR DEVELOPMENT [0001] This invention was made with Government support under Contract No. W911QY-07-C-0032. CROSS REFERENCES [0002] [1] S. R. Nersisyan, N. V. Tabiryan, D. M. Steeves, B. R. Kimball, “Optical Axis Gratings in Liquid Crystals and their use for Polarization insensitive optical switching,” J. Nonlinear Opt. Phys. & Mat., 18, 1-47 (2009). [0003] [2] P. F. McManamon, P. J. Bos, M. J. Escuti, J. Heikenfeld, S. Serati, H. Xie, E. A. Watson, A Review of Phased Array Steering for Narrow-Band Electrooptical Systems, Proceedings of the IEEE, 97, pp. 1078-1096, 2009. [0004] [3] S. R. Nersisyan, N. V. Tabiryan, L. Hoke, D. M. Steeves, B. Kimball, Polarization insensitive imaging through polarization gratings, Optics Express, 17 (3), 1817-1830, 2009. [0005] [4] C. Oh, J. Kim, J. F. Muth, M. Escuti, A new beam steering concept: Riesley gratings, Proc. SPIE, vol. 7466, pp. 74660J1-J8, 2009. [0006] [5] J. C. Wyant, Rotating diffraction grating laser beam scanner, Applied Optics, 14, 1057-1058, 1975. [0000] U.S. Patent Documents 7,319,566 January 2008 Prince et al. 7,324,286 January 2008 Glebov et al. 6,792,028 September 2004 Cook et al. 3,721,486 March 1973 Bramley RIGHTS OF THE GOVERNMENT [0007] The invention described herein may be manufactured and used by or for the Government of the United States for all governmental purposes without the payment of any royalty. FIELD OF THE INVENTION [0008] This invention relates to optical beam control and, in particular, to methods, systems, apparatus and devices for manipulating with light beams, including laser beams and beams with wide spectra and divergence angles, by translating them in the lateral direction and varying their propagation direction over large angles for optical switching, beam scanning, spectral modulation, optical tweezers, thermal seeker, imaging, information displays, and other photonics applications. BACKGROUND OF THE INVENTION [0009] The present invention relates to optical systems for controlling with propagation of light beams. Pointing and positioning systems are enabling components for most laser applications. Conventionally, this is accomplished using mirrors, scan wheels, optical wedges, and other two-axis gimbal arrangements as exemplified, for example, in the U.S. Pat. No. 7,319,566 to Prince et al. These opto-mechanical systems are complex, bulky and heavy for large area beams. For example, the prism apex angle, hence its thickness is increased to achieve larger deflection angles. The electromechanical systems for rotation, translation or oscillation of such mirrors, prisms, and other optical components require high electrical power for their operation. They are relatively slow and have limited range of angles that could be covered within given time period. [0010] Thus, there is a need for thin, light-weight, fast, and inexpensive pointing, positioning, and switching systems for light beams, particularly, for laser beams. The state-of-the-art developments include all-electronics systems and rotating diffraction gratings. The all-electronics systems with no moving parts, as reviewed in P. F. McManamon, P. J. Bos, M. J. Escuti, J. Heikenfeld, S. Serati, H. Xie, E. A. Watson, A Review of Phased Array Steering for Narrow-Band Electrooptical Systems, Proceedings of the IEEE, Vol. 97, pages 1078-1096 (2009), require a large number of high efficiency diffraction gratings and spatial light modulators and/or electrically controlled waveplates. As a result, the overall transmission of these systems is reduced along with their radiation damage threshold, and their speed is limited by the liquid crystal spatial light modulators and variable retarders. [0011] Rotating diffraction gratings as described in J. C. Wyant, “Rotating diffraction grating laser beam scanner,” Applied Optics, 14, pages 1057-1058 (1975), and in the U.S. Pat. No. 3,721,486 to Bramley, partially solves the problem of obtaining larger diffraction angle in thinner optical system, compared, for example to the system of Risley prisms. The light beam diffracted by the first grating in the path of the beam is further diffracted by the second grating. Depending on orientation of those gratings with respect to each other, the deflection angle of the beam can thus be varied between nearly 0 to double of the diffraction angle exhibited by a single grating. The problem with such systems is that phase gratings typically diffract light into multiple orders that need to be blocked along with the order beam. High efficiency Bragg type gratings have narrow spectral and angular range as described in the U.S. Pat. No. 7,324,286 to Glebov et al., and can be used practically for laser beams only, expanded and collimated to minimize divergence. Blazed gratings such as proposed in the U.S. Pat. No. 6,792,028 to Cook et al., still exhibit a multitude of diffraction orders due to their discontinuous structure and do not improve considerably on angular selectivity and efficiency. [0012] The cycloidal diffractive waveplates (DWs), essentially, anisotropic plates meeting half-wave condition but with optical axis orientation rotating in the plane of the waveplate in a cycloidal manner, as described in the review S. R. Nersisyan, N. V. Tabiryan, D. M. Steeves, B. R. Kimball, “Optical Axis Gratings in Liquid Crystals and their use for Polarization insensitive optical switching,” J. Nonlinear Opt. Phys. & Mat., 18, 1-47 (2009), do not have the disadvantages of conventional phase gratings. Moreover, DWs, referred to also as optical axis gratings and polarization gratings, can provide nearly 100% diffraction efficiency in micrometer thin layers. Furthermore, due to their waveplate nature, their diffraction spectrum is broadband, and can even be made practically achromatic. Due to their thinness and high transparency, they can be used in high power laser systems. [0013] Thus, replacing Risley prisms, wedges, mirrors and/or phase gratings with DWs, provides many advantages for manipulating with light beams and imaging. As shown in S. R. Nersisyan, N. V. Tabiryan, L. Hoke, D. M. Steeves, B. Kimball, Polarization insensitive imaging through polarization gratings, Optics Express, 17, 1817-1830 (2009), not only laser beams, but complex images can be steered over large angles without light attenuation or image deformation. That paper further showed that utilizing a pair of closely spaced DWs, one of them with switchable characteristics, it is possible to manipulate with transmission of unpolarized beams and images. This concept suggested and demonstrated in S. R. Nersisyan, N. V. Tabiryan, L. Hoke, D. M. Steeves, B. Kimball, “Polarization insensitive imaging through polarization gratings,” Optics Express, 17, 1817-1830 (2009) was subsequently cited and tested in C. Oh, J. Kim, J. F. Muth, M. Escuti, “A new beam steering concept: Riesley gratings,” Proc. SPIE, vol. 7466, pp. 74660J1-J8 (2009). BRIEF SUMMARY OF THE INVENTION [0014] Thus, the objective of the present invention is providing means for switching and manipulating with light beams and images in lateral and angular space using a set of DWs capable of deflecting nearly 100% of light using thin material layers for a broad band of wavelengths and divergence angles. [0015] The second objective of the present invention is incorporating in said set DWs with controlled characteristics of their optical properties for further enhancing optical manipulation capabilities of said systems. [0016] A further objective of the present invention is providing optical systems, incorporating said DW set, wherein manipulation of light and images with the DW set is transformed into transmission modulation of at the output of the optical system. BRIEF DESCRIPTION OF THE SEVERAL VIEWS OF THE DRAWINGS [0017] FIG. 1A schematically shows deflection of a circularly polarized light beam with a pair of diffractive waveplates. [0018] FIG. 1B schematically shows the structure of diffractive waveplates at different rotational positions. [0019] FIG. 2A shows sample dependence of the propagation angle of a light beam at the output of a pair of diffractive waveplates as a function of the rotational position between the waveplates. [0020] FIG. 2B demonstrates the capability of a pair of diffractive waveplates to steer with no distortions complex images carried by an unpolarized light. [0021] FIG. 3A schematically shows the displacement of a light beam by a pair of diffractive waveplates with parallel orientation of their optical axis modulation directions. [0022] FIG. 3B schematically shows the increase in the resultant deflection angle of a light beam by a pair of diffractive waveplates with anti-parallel orientation of their optical axis modulation directions. [0023] FIG. 3C shows the optical axis orientation pattern in diffractive waveplates with anti-parallel orientation of their optical axis modulation directions. [0024] FIG. 4A schematically shows increasing of the deflection angle of a light beam by a set of four diffractive waveplates each arranged anti-parallel with respect to the previous one. [0025] FIG. 4B demonstrates increasing deflection angle of a light beam by increasing the number of diffractive waveplates from one to four, and comparing them to the original propagation direction of the beam. [0026] FIG. 5 shows increasing deflection angle of a light beam by a system of diffractive waveplates tilted with respect to each other. [0027] FIGS. 6A and B schematically show switching between transmittive and deflective states of a pair of diffractive waveplates when switching one of the diffractive waveplates into an optically homogeneous non-diffractive state shown in C. [0028] FIG. 7 shows a schematic of a beam combining function of a pair of diffractive waveplates. [0029] FIGS. 8A , B and C show a schematic of a system for controlling the spectrum of a light beam with the aid of a set of diffractive waveplates. DETAILED DESCRIPTION OF THE INVENTION [0030] Before explaining the disclosed embodiment of the present invention in detail it is to be understood that the invention is not limited in its application to the details of the particular arrangement shown since the invention is capable of other embodiments. Also, the terminology used herein is for the purpose of description and not limitation. [0031] The preferred embodiment of the present invention includes two DWs, marked with numerals 103 and 105 in FIG. 1A , arranged parallel to each other in close proximity. At the output of the system of DWs 103 and 105 , the pointing direction of the light beam 108 , circularly polarized as shown by spirals 102 and 107 , is, in general, different from that of the propagation direction of the light beam 101 incident on the system, controlled with relative rotational positions of the DWs as schematically shown by arrows 104 and 106 . The optical axis orientation pattern corresponding to different rotational positions of said DWs is shown in FIG. 1B wherein the axes of elongated ellipses 109 correspond to local optical axis orientation direction. In the preferred embodiment, DWs are made of liquid crystal polymers though other optically anisotropic materials and material structures such as subwavelength gratings can be used as well. In general, the layer of DW, typically only a few micrometer thick, is coated on a substrate 110 for stability and robustness. The substrate can be made of a material adequate for the particular application. As an example, a fused silica can be used when controlling UV light beams, and highly transparent glass materials with low absorption can be used for controlling high power laser beams. [0032] The plot of output angles measured for a sample system as a function of angular position between the DWs in S. R. Nersisyan, N. V. Tabiryan, L. Hoke, D. M. Steeves, B. Kimball, “Polarization insensitive imaging through polarization gratings,” Optics Express, 17 (3), 1817-1830 (2009) is shown in FIG. 2A for normal incidence of the beam on the first DW. In the setup shown in FIG. 1A , the polarization of the incident beam is assumed circular, as schematically shown by the spiral 102 . The output beam 108 in this case maintains the circular polarization state 107 . In case of incident unpolarized or linearly polarized beam, two beams of orthogonal circular polarization are generated at the output of the system of two DWs, and the angle between them changes from nearly 0 to nearly double of the diffraction angle depending on relative rotational positions between the DWs as shown in FIG. 2B for light beam carrying a complex image. No image distortions occurs in this process. [0033] Increasing the distance Δz between two identical DWs 302 and 304 , FIG. 3 , introduces transverse shift Δx of the beam 305 emerging from the system with respect to the position of the input beam 301 as a result of deflection of the beam by the first DW 302 . Said emerging beam 305 propagates parallel to the input beam 301 in case the optical axis modulation directions of DWs 302 and 304 are parallel, FIG. 3A , and it also changes in propagation direction when the DW 304 is rotated with respect to DW 302 into a new position 306 , FIG. 3B . The overall deflection angle of the beam can be maximized positioning the output DW 306 anti-parallel with respect to the input DW 302 . The optical axis alignment patterns for anti-parallel DWs 302 and 306 are schematically shown in FIG. 3C . The beam can be steered over arbitrarily large angles by adding DWs into the system. Four DWs, 406 - 409 , are shown in FIG. 4A as an example. The input light 401 undergoes four deflections, 402 - 405 . In order for each subsequent deflection to further increase the resultant deflection angle, the DWs 407 and 409 have to be arranged anti-parallel to DWs 406 and 408 . A demonstration of light deflection by such a system of four DWs is shown in FIG. 4B . In general, DWs can be tilted with respect to each other such as each of the subsequent DWs is nearly perpendicular to the beam deflected by the previous DW. The DWs 507 and 509 are anti-parallel to the DWs 506 and 508 , and all four deflected beam 502 - 505 of the input beam 501 result in increasing total deflection angle. [0034] In another embodiment, one or more DWs in a system can be switched between diffractive and non-diffractive states, optically, thermally, electrically, mechanically, or by any other means, due the effect of external stimuli on optical anisotropy and optical axis orientation modulation pattern. For example, the DW can be made of azobenzene liquid crystal polymer that can be transformed into isotropic state or realigned by light beams as discussed in S. R. Nersisyan, N. V. Tabiryan, D. M. Steeves, B. R. Kimball, “Optical Axis Gratings in Liquid Crystals and their use for Polarization insensitive optical switching,” J. Nonlinear Opt. Phys. & Mat., 18, 1-47 (2009). Alternatively, DWs can be transformed into homogeneous orientation state by electrical fields if they are made of liquid crystals or liquid crystal polymer network materials. [0035] Particularly important is the case shown in FIG. 6 when a DW 603 with a fixed diffractive property is paired with a controllable DW 602 in configuration when their optical axis modulation directions are parallel. As noted in S. R. Nersisyan, N. V. Tabiryan, L. Hoke, D. M. Steeves, B. Kimball, “Polarization insensitive imaging through polarization gratings,” Optics Express, 17 (3), 1817-1830 (2009), this state corresponds to total cancellation of diffraction, and such a pair allows transmitting the light beam 601 through the system as shown schematically in FIG. 6A . An image sensor 604 furnished with an aperture 605 large enough not to block the transmitted beam would not register any distortions to the beam. In case the DW 602 is transformed into a non-diffractive state 606 , the diffraction of, generally, an unpolarized light on the remaining DW 603 redirects the input beam 601 into diffracted beams 607 and 608 as shown in FIG. 6B , diffracting it into orthogonal circular polarized components in case of unpolarized or linearly polarized incident beam. No beam is acting on an image sensor 604 in this case provided the deflected beams propagate beyond the receiving aperture of the image sensor. Thus, the system described in FIG. 6 undergoes switching from high transmission to no or low-transmission state as a result of switching the structure of one of the DWs in the system from diffractive state 603 into a non-diffractive state 606 , FIG. 6C . Indeed, such change in transmission through particular aperture can be obtained also by mechanically changing the rotational position of the DWs or the distance between them. [0036] Paired DWs and their systems can have many applications in photonics. A setup for beam combining is shown in FIG. 7 . Two parallel propagating light beams of orthogonal circular polarizations 701 and 704 , after being deflected by the first DW 707 are further deflected into beams 702 and 705 , emerging as overlapping beams of the same propagation direction 703 and 706 by the second DW 707 in FIG. 7 . [0037] Given the thinness of individual DW layers, a multilayer system can be designed for spectrally selective switching without compromising the high throughput and the small size of the system. In the embodiment shown in FIG. 8 , a set of DW pairs is used for controlling with the spectral content of the transmitted light by allowing light at different portions of the spectrum at least partially be deflected out of the system. The beams 801 and 804 in FIG. 8 are assumed to possess with different, non-overlapping, spectral content. The individual DWs in the first pair 807 are optimized for diffracting the light beam 801 while having diffraction spectrum out of the spectral range of the beam 804 . The individual DWs in the second pair 808 are optimized for diffracting the light beam 804 while having diffraction spectrum out of the spectral range of the beam 801 . Thus, when DWs in both pairs are parallel aligned with respect to their optical axis modulation direction, all the light is transmitted, and the spectral content of the output light is the same as in the input light. In this case shown in FIG. 8A , the input light 801 propagates through the first DW pair into the beam 802 without changing its propagation direction due to diffraction on both DWs constituting the pair 807 . The beam 802 further propagates through the second DW pair 808 into the beam 803 without deflection since its spectrum is out of the diffraction spectrum of the second DW pair 808 . Similarly, the input light 804 propagates through the first DW pair into the beam 805 without changing its propagation direction since its spectrum is out of the diffraction spectrum of the first DW pair 807 . The beam 805 further propagates through the second DW pair 808 into the beam 806 due to the diffraction on both DWs constituting the pair 808 . [0038] In case one of the DWs constituting the first pair 807 is switched into non-diffractive state 809 , or is rotated to double the diffraction angle of the beam 801 by the first DW in the pair 807 , the beam 801 is diffracted out of the optical system into a beam 810 . Propagation of the beam 804 is not affected by that. Thus the light spectrum obtained at the output of the optical system coincides with that of the beam 804 , FIG. 8B . [0039] In case one of the DWs constituting the second pair 808 is switched into non-diffractive state 811 , or is rotated to double the diffraction angle of the beam 805 by the first DW in the pair 808 , the beam 805 is diffracted out of the optical system into a beam 812 . Propagation of the beam 802 is not affected by that. Thus the light spectrum obtained at the output of the optical system coincides with that of the beam 801 , FIG. 8C . [0040] Although the present invention has been described above by way of a preferred embodiment, this embodiment can be modified at will, within the scope of the appended claims, without departing from the spirit and nature of the subject invention.

The objective of the present invention is providing optical systems for controlling with propagation of light beams in lateral and angular space, and through optical apertures. Said light beams include laser beams as well as beams with wide spectrum of wavelengths and large divergence angles. Said optical systems are based on combination of diffractive waveplates with diffractive properties that can be controlled with the aid of external stimuli such as electrical fields, temperature, optical beams and mechanical means.

TECHNICAL FIELD OF THE INVENTION [0001] The present invention lies in the field of molecules known as “small interfering RNA” with therapeutic applications. siRNAs have the ability to reduce gene expression in an extremely specific way (1). These are small sequences of double-strand RNA, normally used in laboratory to modify cell function, which revolutionized cell biology by allowing previously precluded molecular manipulations. STATE OF THE PRIOR ART [0002] CLCN7-dependent ADO2 is a genetic condition affecting 5 individuals in 100,000 newborns (2). It generally appears in teen-aged or adult subjects (3), though various cases of infantile CLCN7-dependent ADO2 are known (4). The disease is characterized by absence of function of bone cells termed osteoclasts (5) and presents with very dense but fragile bones, hemopoiesis and senso-motory function disorders, osteomyelitis and teeth problems. Life expectancy is generally normal; yet, though rarely, early death of the affected individual can occur. Quality of life can instead be markedly compromised due to numerous atraumatic fractures, difficult to reduce surgically, and to an often severely debilitating hematological and neurological symptomatology (6). Cognitive faculties are generally preserved, though significant deficits can appear in the most severe cases (6). CLCN7-dependent ADO2 is characterized by incomplete penetrance, as only about 66% of individuals affected by the mutation manifests the disease (7). It has extremely variable severity, ranging from patients characterized by absence of symptoms to markedly compromised patients (6-8). [0003] CLCN7-dependent ADO2 is due to a mutation of the CLCN7 gene (http://www.ncbi.nlm.nih.gov/nuccore/NM_001114331.2), comprised of 25 hexons and localized in chromosome 16 in humans and in chromosome 17 in mouse, which encodes a protein, termed CIC-7 (http://www.ncbi.nlm.nih.gov/protein/NP_001107803.1), essential to osteoclast function (5), serving for chlorine transmembrane transport. ADO2 is autosomal, as the gene resides in a nonsexual chromosome, and is dominant, as those are point mutations of a gene encoding a homodimeric protein. In point mutations, the entire protein is produced, but it exhibits a change of amino acid which modifies its function. In homodimers there are two identical subunits. Two copies of each gene exist, and in CLCN7-dependent ADO2 only one of the two genes is mutant (mutated), therefore in cells homodimers with both subunits normal, as well as heterodimers with only one subunit mutant, and homodimers with both subunits mutant are formed. Essentially, in affected individuals only one-third of the proteins functions correctly. [0004] In patients in which the mutation determines total absence of the protein, the disease, termed recessive autosomal osteopetrosis, is much more severe because both genes are mutant (9). If instead one of the genes is not mutant, it causes production of normal protein sufficient not to let the disease develop. This condition is known as “haplosufficiency”. [0005] siRNAs (small interfering RNA) are small RNA sequences complementary to specific sequences of messenger RNA (mRNA), inducing its degradation (1). In a preceding work, the present Inventors conjectured the use of mutation-specific siRNAs to silence the expression of the mutated allele of the protein causing CLCN7-dependent ADO2 (10). Even though some siRNAs demonstrated able to block, to a certain extent, the mutated allele expression, their selectivity for the mutated allele and ability to discriminate between mutated allele and wild-type (WT) allele remains an open challenge. [0006] Scope of the present invention is to provide novel siRNAs optimized for the treatment of CLCN7-dependent ADO2 disease. SUMMARY OF THE INVENTION [0007] The invention subject of the present application is based on the ascertainment, carried out by the present Inventors, that complementarity, even if total, to the sequence of mRNA comprising the point mutation is not by itself sufficient to obtain efficacious and selective siRNAs; that is to say, siRNA effective in silencing the expression of the mutated protein but inactive on the expression of the WT protein. [0008] On the contrary, the optimum combination of efficacy and selectivity depends on plural factors, such as the mutation itself on the mRNA, the length of the sequences flanking the mutation, therefore the position of the mutated nucleotide in the siRNA sequence, the presence or absence of one or more nucleotide mismatches compared to the WT sequence of the mRNA and the position of said mismatch in the siRNA sequence: in short, from the design to the sequence itself of the siRNA. [0009] Therefore, a first object of the present application are small interfering RNA (siRNA) and their derivatives or their precursors complementary to the region comprising a point mutation in the messenger RNA (mRNA) of the mutated human gene CLCN7. The siRNA object of the application are characterized in that (i) said mutations reflect corresponding mutations of the CIC-7 protein: Y99C, D145G, W179X, G203D, L213F, L213F, L213F, G215R, P249L, R286W, R286Q, P470Q, R409W, L490F, G677V, 688del, K689E, R762L, G765B, L766P, R767W, A788D; that (ii) the siRNAs have a nucleotide sequence comprising a fragment of 15 to 25 nucleotides, comprising the point mutation; that (iii) the siRNA selectively reduce the expression of mutated CIC-7 proteins and that (iv) the ratio of efficacy of the siRNA of the invention in reducing the expression of mutated CIC-7 protein compared to the normal protein is greater than one. [0010] Optionally, the sequence of the small interfering RNA (siRNA) of the invention comprises, in addition to the mutated nucleotide, one or more nucleotide mismatches compared to the corresponding target sequence of the mRNA containing the mutation. Optionally, the sequence of the small interfering RNA (siRNA) of the invention also comprises a short sticky sequence to the 3′ end, consisting of nucleotides dA and dT. A second object of the invention is represented by the above-indicated siRNA for use in a therapeutic treatment, specifically in the therapeutic treatment of ADO2 caused by a mutation of the CLCN7 gene. [0011] A third object of the invention is a method for the preparation of the above-disclosed siRNAs. [0012] A fourth object of the invention is represented by pharmaceutical compositions comprising, as active ingredient, one or more siRNAs and a pharmacologically acceptable excipient. Such compositions are preferably for parenteral administration. A further object of the invention is represented by the same compositions for use in the therapeutic treatment of ADO2, also in association with a second active ingredient. [0013] The siRNA according to the invention, designed and tested in isolated cells and in an animal model, proved highly specific for the mutated gene. They selectively eliminate up to 95% of the transcript of the mutated gene, creating a situation, similar to haplosufficiency, which restores osteoclast function and redresses disease symptoms. The siRNA of the invention moreover afford the further advantages of being internalized by osteoclastic cells by mere incubation, with no need of any transfection agent, and of remaining in the cell for a long time. DESCRIPTION OF THE FIGURES [0014] FIG. 1 —Scheme of WT-CLCN7 construct, obtained by cloning of the complete sequence of human CLCN7 cDNA in the pEGFP-C1 expression vector, by restriction enzymes HindIII and XhoI. [0015] FIG. 2 —Construct sequences checked by direct DNA sequencing. Standard sequences are reported in the upper row, the mutated nucleotide is shown in the row below; the numbers preceding and following each sequence indicate the position of the same inside the cDNA derived from mRNA for the human CLCN7 gene (http://www.ncbi.nlm.nih.gov/nuccore/NM_001114331.2). [0016] FIG. 3 —HEK293 cells were transfected with empty vector or with WT-, p.R767W-, p.G215R-, p.R286W- and p.A788D-EGFP vector. CLCN7 relative expression was quantified by real-time RT-PCR. Mean±s.e. normalized for GAPDH *p<0.000001 vs. cells transfected with empty EGFP plasmid (Student's t test). [0017] FIG. 4 —HEK293 cells were transfected with vectors carrying WT-EGFP, p.R767W-EGFP and p.G215R-EGFP constructs, then the expression and the localization of the fluorescent fusion protein EGFP was detected by confocal microscopy along with the expression of Golgi apparatus markers (gamma-adaptin), early endosomes (EEA1), acid vesicles (Lysotracker Red) and lysosomes (lamp-1). Co-localization of the fusion protein EGFP with the indicated markers is shown in the “merge” panels. Objective lens magnification=63X. Similar results were obtained also for mutations p.A788D and p.R286W (not shown). [0018] FIG. 5 —RAW264.7 cells were transfected with empty EGFP vector or with p.R767W-EGFP vector and differentiated into osteoclasts on bovine bone slices, by treatment with 100 ng/ml RANKL. The slices were then sonicated to remove cells and assessed for resorption lacuna (pit) formation (pit index) after staining with 0.1% toluidine blue. Mean±s.e. *p<0.05 vs. cells transfected with empty EGFP vector (Student's t test). [0019] FIG. 6 —Human osteoclasts were differentiated from peripheral blood mononuclear cells of a healthy donor, by incubation for 14 days with 20 ng/ml M-CSF and 30 ng/ml RANKL. Cells were transfected with WT-EGFP vector by the AMAXA method. After 2 days, the fluorescence of the EGFP fusion protein was detected by confocal microscopy. The correct co-localization with the markers indicated for the Golgi apparatus (gamma-adaptin), early endosomes (EEA-1), acid vesicles (Lysotracker Red) and lysosomes (lamp-1) is shown in the “merge” panels. Objective lens magnification=63X. [0020] FIG. 7 —HEK293 cells transfected with WT vectors or carrying the mutation p.R767W of the CLCN7 gene and treated with control scrambled siRNA or with increasing concentrations of p.R767W 1 siRNA for 48 hours. At the end of the incubation, the RNA was extracted and the expression of the CLCN7 transcript was assessed by real-time RT-PCR. Mean±s.e. normalized for GAPDH, expressed as percentage over the treatment with scrambled siRNA (point 0 on the X-axis). *p<0.05 vs. cells transfected with p.R767W- and WT-EGFP treated with control (SCR, scrambled) siRNA (Student's t test). [0021] FIG. 8 —RAW264.7 cells were transfected with empty vector or with WT- or p.R767W-EGFP vectors, differentiated into osteoclasts on bovine bone slices and treated with control (SCR) siRNA, or with siRNA specific for the transcript bearing the mutation p.R767W (p.R767W 1 RNA). Quantification of bone resorption shows the ability of p.R767W 1 siRNA to improve bone resorption reduced by the mutated construct (compare black and light grey bars). Instead, p.R767W 1 siRNA did not significantly modify bone resorption in cells transfected with WT-EGFP vector (compare white and dark grey bars). Mean±s.e. *p<0.05 vs. cells transfected with empty EGFP vector (Student's t test). [0022] FIG. 9 —HEK293 cells, transfected with WT- and p.R767W-EGFP vectors, were treated with control scrambled siRNA or with increasing concentrations of the siRNAs for p.R767W indicated in the figure (sequences in Table 2) and evaluated for CLCN7 transcript expression by real-time RT-PCR. p.R767W siRNA 2C showed greater efficacy and specificity compared with the other p.R767W siRNAs. Mean±s.e. normalized for GAPDH, expressed as percentage over the treatment with scrambled siRNA (point 0 on the X-axis). *p<0.05 vs. cells transfected with p.R767W-EGFP and treated with control (SCR, scrambled) siRNA. # p<0.05 vs. cells transfected with WT-EGFP vector and treated with control (SCR) siRNA (Student's t test). [0023] FIG. 10 —(A) Human osteoclasts were differentiated from peripheral blood mononuclear cells of a healthy donor and incubated with 300 nM Cy3-CLCN7 WT siRNA for 48 hours. Cells were then fixed with 4% paraformaldehyde and Cy3 fluorescence inside osteoclasts was detected by confocal microscopy. Objective lens magnification=63X. (B,C) Human osteoclasts were transfected with empty vector or with p.R767W-EGFP vector. Osteoclasts transfected with p.R767W-EGFP vector were treated for 48 hours with control (SCR) siRNA or p.R767W siRNA 2C (Table 2) at the indicated concentrations. After 48 hours, expression (B) of the EGFP transcript normalized for GAPDH, and (C) bone resorption were analyzed, respectively by real-time RT-PCR and pit assay. Mean±s.e. # p<0.005 vs. osteoclasts transfected with p.R767W-EGFP and treated with control (SCR) siRNA. *p<0.05 vs. osteoclasts transfected with empty EGFP vector (Student's t test). [0024] FIG. 11 —Human osteoclasts were differentiated from peripheral blood mononuclear cells of a healthy donor and treated for 48 hours with the concentrations of WT siRNA or p.R767W siRNA 2C indicated in the figure. RNA was extracted and subjected to RT-PCR. Note the absence of effect of p.R767W siRNA 2C on expression of normal CLCN7 transcript (normalized for GAPDH). [0025] FIG. 12 —Human osteoclasts were differentiated from peripheral blood mononuclear cells of a healthy donor, plated on bone slices, incubated and treated for 48 hours with scrambled siRNA or with the indicated concentrations of p.R767W siRNA 2C. Bone resorption was then quantified by pit assay. Mean±s.e. expressed as percentage over the treatment with scrambled siRNA (point 0 on the X-axis). Note the absence of a statistically significant effect in resorption pit formation by the treatment with p.R767W siRNA 2C (Student's t test). [0026] FIG. 13 —(A) HEK293 cells and (B) primary human osteoclasts were transfected with WT- or p.G215R-EGFP vectors and treated for 48 hours with 100 nM of control (SCR, scrambled) siRNA, or of p.G215R-EGFP (p.G215R)-specific siRNA. CLCN7 mRNA expression was assessed by real-time RT-PCR, using primers specific for CLCN7, for HEK239 cells (which do not express the endogenous gene CLCN7), and for EGFP for human osteoclasts (to distinguish the mutated exogenous CLCN7 gene from the normal endogenous one). Mean±s.e. normalized with GAPDH (Student's t test). [0027] FIG. 14 —Primary osteoclasts, generated from the bone marrow mononuclear cells of WT- and p.G213R-clcn7 ADO2 (p.G213R KI) mice by incubation with 50 ng/ml M-CSF and 120 ng/ml RANKL, were plated on bone slices and treated for 7 days with control (SCR) siRNA, or p.G213R-clcn7 (p.G213R) specific siRNA, at the concentration of 100 nM. At the end of the incubation, cells were removed by sonication, slices were stained with 0.1% toluidine blue and bone resorption was evaluated by pit assay. Mean±s.e. *p=0.0001 vs. WT, #p=0.003 vs. p.G213R KI (Student's t test). [0028] FIG. 15 —1 month-old CD1 mice were treated with control scrambled siRNA or with 2 mg/Kg of body weight of clcn7-siRNA by TransIT-QR Hydrodynamic Delivery Kit. After 24 and 48 hours the animals were sacrificed and RNA was extracted from kidney, brain, liver, lung, heart, spleen and tibia. clcn7 gene transcriptional levels were analyzed by real-time RT-PCR, normalized for gapdh and expressed as percentage over the treatment with scrambled siRNA (point 0 on the X-axis). [0029] FIG. 16 —(A) RT-PCR using primers specific for the p.G213R mRNA (Fw: CAAGTGCTTCCTCAATG (SEQ ID NO:32); Rv: GCCCTCTTCCAAGCTAAA (SEQ ID NO:33) showing transcript amplification only in primary osteoclasts of heterozygous and homozygous p.G213R KI mice, while in wild-type (WT) osteoclasts no transcript appears amplified. (B) Direct DNA sequencing of the amplified transcript shown in figure (A) in heterozygous p.G213R/WT osteoclasts, demonstrating only the mutated sequence. [0030] FIG. 17 —Osteoclasts generated from bone marrow mononuclear cells of WT and p.G213R KI mice were treated with the indicated concentration of control (SCR) siRNA or p.G213R-clcn7-specific siRNA. Real-time RT-PCR was performed using the primers specific for the mutated transcript indicated in FIG. 14 . Mean±s.e. (Student's t test) [0031] FIG. 18 —Three month-old p.G213R KI mice received an intraperitoneal (i.p.) injection of 4 mg/kg of p.G213R-clcn7 sticky siRNA/jetPEI™ (conjugate) and were sacrificed at the indicated time points. Sera were collected and evaluated for total RNA levels by Nanodrop. Mean±s.e. (ANOVA). [0032] FIG. 19 —Ten day-old p.G213R KI mice received an i.p. injection of 4 mg/kg of p.G213R-clcn7 sticky siRNA/jetPEI® conjugate, 3 times a week for 4 weeks. At the end of the experiment mice were sacrificed, RNA was extracted from the organs indicated in figure and subjected to RT-PCR, using primers specific for the mutant transcript indicated in FIG. 14 , normalized with gadph. Mean±s.e. (Student's t test) [0033] FIG. 20 —p.G213R ADO2 mice were subjected to i.p. injection with control (SCR, scrambled) siRNA or p.G213R-clcn7 sticky siRNA/jetPEI (p.G213R), at the doses indicated on the X-axis. After 48 hours, tibias were collected, RNA was extracted and the levels of p.G213R-clcn7 mutated mRNA were evaluated by real-time RT-PCR using the pair of primers specific for the mutated sequence indicated in FIG. 14 . Mean±s.e. normalized for GAPDH. [0034] FIG. 21 —p.G213R KI mice received i.p. injections of 4 mg/kg of SRC-siRNA or p.G213R-clcn7 sticky siRNA/jetPEI, 3 times a week for 4 weeks. At the end of the experiment, mice were sacrificed and subjected to histopathological evaluation of the organs indicated in Figure by haematoxylin/eosin staining (Bar=100 μm for spleen and kidney, 20 μm for liver). [0035] FIG. 22 —Sera were collected from the mice described in FIG. 20 and analyzed by Refloton method for the biomarkers of renal (uric acid) and hepatic [glutamic oxaloacetic transaminase (GOT)] functions, and for the ADO2 biomarker creatine kinase (CK). Normal values are comprised between the dotted lines. Mean±s.e. (Student's t test). [0036] FIG. 23 —WT and p.G213R KI mice received i.p. injections of 4 mg/kg of SRC-siRNA or p.G213R-clcn7 sticky siRNA/jetPEI, 3 times a week for 2 and 4 weeks. At the end of the experiment, mice were sacrificed and sera were collected for evaluating the levels of osteoclastic (isoform 5b of TRAcP enzyme) and bone resorption (CTX) biomarkers and for calculating the CTX/TRAcP ratio. [0037] FIG. 24 —Analysis of the bone phenotype of mice treated for 2 weeks as indicated in FIG. 22 . (A) μCT analysis of the proximal region of the tibia. (B) trabecular bone volume over total tissue volume (BV/TV), (C) Trabecular number (Th.N). (E) Trabecular thickness (Tb.Th). (F) Trabecular separation (Tb.Sp). Mean±s.d. of 4-7 mice/group (Student's t test). [0038] FIG. 25 —Ten day-old WT and p.G213r KI mice received i.p. injections of 4 mg/kg of SRC-siRNA or of p.G213R-clcn7 sticky siRNA/jetPEI, 3 times a week for 4 weeks. At the end of the experiment mice were sacrificed and their bone phenotype analyzed. (A) μCT of the proximal region of the tibia. (B) Trabecular bone volume over total tissue volume (BV/TV). (C) Trabecular number (Tb.N). (D) Trabecular thickness (Tb.Th). (E) Trabecular separation (Tb.Sp). (F) Serum concentration of parathyroid hormone (PTH). [0039] FIG. 26 —Analysis of osteoclastic phenotype in mice described in FIG. 24 . (A) Histochemical TRAcP enzyme staining to highlight osteoclasts (purple cells). Bar=100 μm. (B) Osteoclast surface over bone surface (Oc.S/BS). (C) Osteoclast number over bone perimeter (Oc.N/B Pm). (D) Transcriptional expression, by real-time RT-PCR on RNA extracted from the femurs, of osteoclast [Tracp and Cathepsin K (CatK)] and osteoblast [Alkaline phosphatase (ALP) and Runt-related transcription factor 2 (Runx 2)] genes normalized with gapdh. (E) Osteoclast-eroded surface over total bone surface (ES/BS). Mean±s.d. (Student's t test). [0040] FIG. 27 —Analysis of cortical, cartilagineous and osteoblastic parameters in mice described in FIG. 24 . (A) Cortical bone thickness (Cor.Th). (B) Growth plate thickness (width). (C) Osteoblast surface over bone surface (Ob.S/BS). (D) Histological images of osteoid (arrows). Bar=5 μm. (E) Osteoid volume over bone volume (OV/BV). (F) Calcein labeling (green fluorescence) of mineral deposition (double arrowheads). Bar=2 μm. (G) Mineral apposition rate (MAR). (H) Mineralized surface over bone surface (MS/BS). (I) Bone formation rate (BFR). [0041] FIG. 28 —Analysis of bone quality by indentation in mice described in FIG. 25 . (A) Total indentation distance (TID). (B) First cycle indentation distance (ID). (C) Touchdown distance (TDD). Mean±s.d. of 3-7 mice/group. (Student's t test). [0042] FIG. 29 —Human osteoclasts were transfected with the expression vectors indicated in Figure and treated for 48 hours with the concentrations, indicated on the X-axis, of (A) p.G213R-, (B) p.R767R- and (C) p.R286W-EGFP-specific siRNAs. Real-time RT-PCR using primers specific for EGFP, normalized with GAPDH. (D) Osteoclasts were generated from blood mononuclear cells of an ADO2 patient carrying the p.G215R mutation, cultured on bovine bone slices and treated with the indicated concentration of SCR-siRNA and p.G215R-siRNA. At the end of the experiment, cells were removed by sonication and bone resorption evaluated by measuring the resorption pits. Results of a single experiment without replicates. [0043] FIG. 30 —(A) Human breast cancer cells MDA-MB-231 (MDA) were transfected with WT- or p.R767W-EGFP vectors and treated for 48 hours with control scrambled siRNA or with p.R767W siRNA 2C at the concentrations indicated on the X-axis. The graph shows the results of the expression of EGFP conjugated to transfected CLCN7 gene, analyzed by real-time RT-PCR performed using a pair of primers specific for EGFP normalized for GAPDH. Mean±s.e. expressed as percentage over treatment with scrambled siRNA (point 0 on the X-axis). p=0.02 (statistical test: area below curve). (B) Xenotumors were obtained by subcutaneous injection of human breast cancer cells, MDA-MB-231, transfected with p.R767W-EGFP vector, in the sides of Balb/c nu/nu athymic (immunocompromised) mice. When tumors reached the volume of 1 cm 3 , the mice were treated once for 96 hours with vehicle (jetPEI) or 4 mg/Kg body weight of control (SCR, scrambled) siRNA, WT-siRNA or R767W siRNA 2C complexed with jetPEI. Real-time RT-PCRs were then performed on RNA extracted from the tumors, using a pair of primers specific for EGFP. Mean±s.e. normalized for GAPDH. p=0.02 (Student's t test). DETAILED DESCRIPTION OF THE INVENTION [0044] It is known that the mRNA of human CLCN7 gene can comprise mutations, pathogenetic ones in ADO2, which generate mutated proteins as indicated in Table 1 [0000] TABLE 1 Mutation and position in Amino acid the protein substitution Other mutation 1 Y99C Tyrosine/Cysteine 2 D145G Aspartic acid/Glycine 3 W179X Tryptophan/Unknown 4 G203D Glycine/Aspartic acid 5 L213F Leucine/Phenylalanine 6 G215R Glycine/Arginine 7 P249L Proline/Leucine 8 R286W Arginine/Tryptophan 10 R286Q Arginine/Glutamine 11 P470Q Proline/Glutamine 12 R409W Arginine/Tryptophan 13 L490F Leucine/Phenylalanine 14 G677V Glycine/Valine 15 688del — Amino acid deletion in position 688 of the protein 16 K689E Lysine/Glutamic acid 17 R762L Arginine/Leucine 18 G765B Glycine/ 19 L766P Leucine/Proline 20 R767W Arginine/Tryptophan 21 A788D Alanine/Aspartic acid 22 2423delAG — Adenine/Guanine deletion in position 2423 of the DNA [0045] Other potential mutations of the CLCN7 gene that might result into possible muteins from the CIC-7 protein and in as many siRNAs according to the invention are the following ones: R223L, R223P, R223G, R223K, R223W, R223I, R223M, R223C, R223S, R265L, R265P, R265G, R265K, R265W, R265I, R265M, R265C, R265S, R271L, R271P, R271G, R271K, R271W, R271I, R271M, R271C, R271S, R280L, R280P, R280G, R280K, R280W, R280I, R280M, R280C, R280S, R281L, R281P, R281G, R281K, R281W, R281I, R281M, R281C, R281S, R286L, R286P, R286G, R286K, R286I, R286M, R286C, R286S, R326L, R326P, R326G, R326K, R326W, R326I, R326M, R326C, R326S, R362L, R362P, R362G, R362K, R362W, R362I, R362M, R362C, R361S, R403L, R403P, R403G, R403K, R403W, R403I, R403M, R403C, R403S, R405L, R405P, R405G, R405K, R405W, R405I, R405M, R405C, R405S, R409L, R409P, R409G, R409K, R409W, R409I, R409M, R409C, R409S, R436L, R436P, R436G, R436K, R436W, R436I, R436M, R436C, R436S, R526L, R526P, R526G, R526K, R526W, R526I, R526M, R526C, R526S, C211F, C211S, C211Y, C211R, C211G, C211W, C411F, C411S, C411Y, C411R, C411G, C411W, C438F, C438S, C438Y, C438R, C438G, C438W, W541R, W541S, W541L, W541G, W616R, W616S, W616L, W616G, L224S, L224P, L224W, L224H, L224Q, L224R, L224F, L224I, L224M, L224V, L224S, L224P, L224W, L224H, L224Q, L224R, L224F, L224I, L224M, L224V, L227S, L227P, L227W, L227H, L227Q, L227R, L227F, L227I, L227M, L227V, L564S, L564P, L564W, L564H, L564Q, L564R, L564F, L564I, L564M, L564V, S290Y, S290C, S290W, S290F, S290P, S290L, S290T, S290A, S290N, S365Y, S365C, S365W, S365F, S365P, S365L, S365T, S365A, S365N, S473Y, S473C, S473W, S473F, S473P, S473L, S473T, S473A, S473N, G241R, G241S, G241W, G241C, G241D, G241E, G241A, G241V, G347R, G347S, G347W, G347C, G347D, G347E, G347A, G347V, G361R, G361S, G361W, G361C, G361D, G361E, G361A, G361V [0046] siRNAs [0047] Small interfering RNAs complementary to the region comprising a point mutation in the messenger RNA (mRNA) of human CLCN7 gene were designed and produced for all gene mutations, known to be pathogenetic for ADO2, reported in Table 1. [0048] The small interfering RNA (siRNA) of the invention are double-strand (duplex) sequences, of which the first one is termed “guide” (or antisense) and the second one “passenger” (or sense). The guide strand (antisense) is that complementary to the target RNA that is to be inhibited, silenced or degraded. [0049] As the sequence of the passenger strand is complementary to the guide strand, for all siRNAs of the invention indicated in the present application, only the sequence of the guide strand is reported. The siRNAs of the invention have a sequence comprising or consisting in a fragment composed of 15 to 29 nucleotides, e.g. 16, 17, 18, 19, 20, 21, 22, 23 or 24, 25, 26, 27 or 28 containing the point mutation; [0050] The siRNAs of the invention are selected for their ability to selectively bind to the mRNA transcribed from the mutated allele forms of the CLCN7 gene, reducing or suppressing the expression of mutated protein CIC-7. Thanks to their selectivity of silencing of the mutated gene, their efficacy in the reduction of the expression is greater for the mutated protein than for the normal protein. Therefore, they exhibit a mutated CIC-7/normal CIC-7 ratio of efficacy greater than one. [0051] In order to further enhance siRNAs selectivity and/or specificity to mutated mRNA, siRNAs sequence can comprise one or more nucleotides non-complementary (mismatch) to said mutated RNA sequence. With this strategy, novel siRNAs were designed (Table 2). Various siRNA exhibited significantly greater specificity for mutated mRNA, compared to W.T. mRNA. With one of them (termed p.R767W siRNA 2C), an efficacy of 90% reduction of mutated mRNA was obtained, without any reduction of WT mRNA ( FIG. 9 ). [0052] siRNA Derivatives [0053] Moreover, in order to increase the stability of the same siRNAs produced and to improve the efficiency of the produced effect, one or more nucleotides forming the siRNAs sequence can be chemically modified in order to obtain derivatives of the siRNA of the invention. All derivatives described hereinafter are therefore encompassed by the protective scope of the present application. [0054] Firstly, the siRNA sequence can be provided with a dTdT or dAdT sequence protruding to the 3′ end. The latter sequence, besides lending stability and improving the efficiency, induces siRNA oligomerization in order to mimick the DNA (sticky siRNA). Sticky siRNAs can therefore be associated with usual reagents ensuring efficient siRNA distribution in vivo and decreasing the ability to cause immune responses mediated by pro-inflammatory cytokines and interferon: for instance, the jetPEI® product, which is a linear polyethylenimine derivative provided by PolyPlus Transfection. [0055] In the in vivo assays performed on a murine ADO2 model, just the sticky siRNA/jetPEI conjugates were used. However, the siRNAs of the invention, non-modified or differently modified, as described hereinafter, can equally be used efficaciously. [0056] Other derivatives improving the stability of the siRNAs of the invention in the form of duplex are the 2′-alcoxy (C1, C2, C3, C4) derivatives, e.g. the 2′-methoxy-derivatives, (i.e. 2′-OMe derivatives) (Denise M Kenski, Gabor Butora, Aarron T Willingham, Abby J Cooper, Wenlang Fu, Ning Qi, Ferdie Soriano, Ian W Davies and W Michael Flanagan. “siRNA-optimized Modifications for Enhanced In Vivo Activity.” Molecular Therapy Nucleic Acids (2012) 1, e5; doi:10.1038/mtna.2011.4). 2′-OMe-derivatives, normally present in rRNA and in tRNA, are atoxic derivatives of the siRNA of the invention, wherein the —OMe group is inserted in position 2′ of the ribose core in the sense- or antisense strand, or in both. [0057] Also 2′-fluorine (i.e. 2′-F) -derivatives (Denise M. Kenski et al, supra) are compatible with the function carried out by the siRNA of the invention and enhance the stability of the duplex thereof against nuclease degradation. Fluorine incorporation in position 2′ of the ribose core maintains siRNAs activity both in vitro and in vivo, increasing their stability. Combined use of 2′-F in pyrimidine nucleotides with 2′-OMe in purine nucleotides results in a duplex siRNA of extreme in-serum stability and markedly improved efficacy. [0058] 2′-O-(2-methoxyethyl) RNA derivatives (MOE-RNA) (Mark A. Behlke. “Chemical Modification of siRNAs for In Vivo Use”. Oligonucleotides 18:305-320 (2008)) can equally be used to enhance the stability of the siRNA of the invention. MOE groups are frequently used in antisense oligonucleotides to give to the oligonucleotide high resistance to nucleases and to increase Tm. [0059] Other siRNA derivatives, having improved function and stability, suitable to the present invention, are the 2′-O-benzyl derivatives and the 2′-O-methyl-4-pyridine (see Denise M. Kenski et al supra), 2′-amino (2′-NH), 2′-aminoethyl (2′-AE), 2′-guanidinopropyl(2′-GP) derivatives. [0060] Particularly interesting to the ends of the present invention, due to their stability, are the LNAs (locked nucleic acids) derivatives of siRNA (see Mark A. Behlke, supra). As well-known to a person skilled in the art, these derivatives are characterized by a methylene bridge between ribose positions 2′-O and 4′-C. The methylene bridge blocks the saccharide unit into the 3′-endo configuration, thereby affording a significant Tm increase and resistance to nucleases. [0061] Precursors [0062] In a specific embodiment of the invention, the siRNAs or derivatives thereof can be used in the form of their precursors in vivo. The latter are also an object of the present invention. [0063] By way of example, siRNAs can be replaced by the corresponding shRNA (short hairpin RNA), in particular in the scope of gene therapy. As is well-known to a person skilled in the art, shRNAs are short RNA sequences or transcripts, consisting in a double-strand structure formed by the coupling of two complementary sequences of about 15-29 nucleotides each, normally 19-25 or 15-20, linked by a loop of about 2-10 nucleotides, e.g. 4-9 or 5-6 nucleotides. When introduced and expressed into the cell, the shRNA-forming transcripts are processed by the enzymatic complex DICER, which by cutting the loop sequence converts, directly into the cell, the shRNAs into the corresponding siRNAs. The latter will then carry out their target gene silencing or knockdown function. Therefore, within the scope of gene therapy, the siRNAs of the invention can be replaced by the corresponding shRNAs. [0064] All of the above-described derivatives and precursors are encompassed within the protective scope of the present application. [0065] Within the scope of the present work, various siRNA specific for the mRNA of alleles of human CLCN7 gene containing the mutations indicated in Table 1, or the murine gene mutation p.G213R, were designed and produced. Then, the efficacy of the individual RNA fragments reported hereinafter in inhibiting the expression of protein CLC7-WT compared to proteins mutated in positions: p.G215R, p.R767W, p.R286W, p.A788D, was analyzed. [0000] (SEQ ID NO: 1) UUCCUCAA U A GGGUGAAGA (SEQ ID NO: 2) UUCCUCAA A GGGUG AGA (SEQ ID NO: 3) UUCCUCAA A GGGU AAGA (SEQ ID NO: 4) UUCCUCAA A GGUGAAG (SEQ ID NO: 5) UUCCUCAA A G GUGAAGA (SEQ ID NO: 6) UUCCUCAA A GGGUGA GA (SEQ ID NO: 7) UUCCUCAAC A GGGUGAA A (SEQ ID NO: 8) CAAC A G GUGAAGAUCCCC (SEQ ID NO: 9) UUCCUCAAC A GGGUGAAGA (SEQ ID NO: 10) CUCAAC A GGGUGAAGAUCC (SEQ ID NO: 11) CAAC A GGGUGAAGAUCCCC (SEQ ID NO: 12) AAC A GG UGAAGAUCCCCC (SEQ ID NO: 13) AAC A GGGUGAAGAUCCCCC (SEQ ID NO: 14) CCUGGGCCUG U GGCACCUG (SEQ ID NO: 15) CCUGGGCCUG U GGCACCU (SEQ ID NO: 16) CCUGGGCCUG U GGC CCUG (SEQ ID NO: 17) CCUGGGCCUG U GGCA CUG (SEQ ID NO: 18) ACAGAGAAG U GGGACUUCG (SEQ ID NO: 19) ACAGAGAAG U GGGACUUC (SEQ ID NO: 20) ACAGAGAAG U GGG GCUUCG (SEQ ID NO: 21) ACAGAGAAG U GGGA UUCG (SEQ ID NO: 22) AGGACCUCG A CAGGUACCG (SEQ ID NO: 23) AGGACCUCG A CAGGUACC (SEQ ID NO: 24) AGGACCUCG A CAG UACCG (SEQ ID NO: 25) AGGACCUCG A CAGG ACCG (SEQ ID NO: 26) AGGACCUCG A C GGUACCG (SEQ ID NO: 27) AGGACCUCG A CAGGUA CG (SEQ ID NO: 28) AGGACCUCG A CAGGU CCG (SEQ ID NO: 29) GGA CUCG A CAGGUACCGC [0066] In the experimental work described in the present application, all siRNA sequences reported above were equipped with a dTdT sequence protruding to the 3′ end to improve their stability and efficacy. For in vivo use, the dTdT sequence is replaced by the dAdT sequence, which further improves the stability and efficacy thereof and enables the binding thereof to any vehicle allowing an improved in vivo distribution of the siRNA and reducing any immune response. [0067] The results obtained with the siRNAs assayed in vitro, in terms of their efficacy on WT mRNA or on mutated mRNA, are listed in Table 2 below. [0000] TABLE 2 Number of mis- matches Efficacy Efficacy siRNA vs vs wild- vs Mutation name wild-type 5′ sense-sequence type Mutant p.G215R G215R 2 UUCCUCAA A GGGUGAAGA- + 0 1M dTdT 3′ G215R 3 UUCCUCAA A GGGUG AGA- + 0 2A dTdT 3′ G215R 3 UUCCUCAA A GGGU AAGA- ++ 0 2B dTdT 3′ G215R 3 UUCCUCAA A GGUGAAG - +++ 0 2C dTdT 3′ G215R 3 UUCCUCAA A G GUGAAGA- 0 ND* 2D dTdT 3′ G215R 3 UUCCUCAA A GGGUGA GA- 0 ND* 2E dTdT 3′ G215R 2 UUCCUCAAC A GGGUGAA A- ++ + 2F dTdT 3′ G215R 2 CAAC A G GUGAAGAUCCCC- ++ ++ 2G dTdT 3′ G215R 1 UUCCUCAAC A GGGUGAAGA- + + 2H dTdT 3′ G215R 1 CUCAAC A GGGUGAAGAUCC- ++ + 2I dTdT 3′ G215R 1 CAAC A GGGUGAAGAUCCCC- + ++ 2L dTdT 3′ G215R 2 AAC A GG UGAAGAUCCCCC- ++ + 2N dTdT 3′ G215R 1 AAC A GGGUGAAGAUCCCCC- 0 +++ 2M dTdT 3′ p.R767W R767W 1 1 CCUGGGCCUG U GGCACCUG- ++ ++ dTdT 3′ R767W 2 CCUGGGCCUG U GGCACCU - + ++ 2A dTdT 3′ R767W 2 CCUGGGCCUG U GGC CCUG- + +++ 2B dTdT 3′ R767W 2 CCUGGGCCUG U GGCA CUG- 0 ++++ 2C dTdT 3′ p.R286W R286W 1 1 ACAGAGAAG U GGGACUUCG- +++ ++ dTdT 3′ R286W 2 ACAGAGAAG U GGGACUUC - + ++ 2A dTdT 3′ R286W 2 ACAGAGAAG U GGG CUUCG- 0 ++ 2B dTdT 3′ R286W 2 ACAGAGAAG U GGGA UUCG- + ++++ 2C dTdT 3′ p.A788D A788D 1 1 AGGACCUCG A CAGGUACCG- ++++ +++ dTdT 3′ A788D 2 AGGACCUCG A CAGGUACC - ++++ +++ 2A dTdT 3′ A788D 2 AGGACCUCG A CAG UACCG- ++++ +++ 2B dTdT 3′ A788D 2 AGGACCUCG A CAGG ACCG- ++++ +++ 2C dTdT 3′ A788D 2 AGGACCUCG A C GGUACCG- ++++ ++ 2D dTdT 3′ A788D 2 AGGACCUCG A CAGGUA CG- ++++ ++ 2E dTdT 3′ A788D 2 AGGACCUCG A CAGGU CCG- 0 0 2F dTdT 3′ A788D 2 GGA CUCG A CAGGUACCGC- 0 + 2G dTdT 3′ [0068] The sequences of siRNAs used in the in vitro study report in bold underlined the mutant nucleotide (referred to human sequence http://www.ncbi.nlm.nih.gov/nuccore/NM_001114331.2) and in bold italics the additional mismatch nucleotide/s. siRNAs indicated in bold in the second column are those deemed most effective and specific. The assays were performed in HEK293 human cells, apart from siRNAs marked with (*), assayed exclusively in human WT osteoclasts, whose efficacy on the corresponding transcript was not determined (ND). Also a siRNA recognizing the mRNA of CLCN7 gene (http://www.ncbi.nlm.nih.gov/nuccore/NM_011930.3) was produced, mutated in position 213 of mouse protein (p.G213R) (corresponding to mutation p.G215R in humans) as indicated in Table 3. [0000] TABLE 3 Number of siRNA mismatches siRNA name vs wild-type 5′ sense sequence p.G213 p.G213R- 2 AA A GGGUGAAGAUCCCCCdAdT3′ (SEQ R clcn7 ID NO: 30) WT clcn7 0 AAUGGGGUGAAGAUCCCCCdAdT3′ (SEQ ID NO: 31) [0069] siRNA sequence for gene p.G213R-clcn7 and for the normal gene (clcn7) to be used in vivo. Bold underlined=mutant nucleotide; Bold italics=mismatch nucleotide. [0070] The preferred siRNAs of the invention are: [0000] (SEQ ID NO: 11) CAAC AGGGUGAAGAUCCCC (SEQ ID NO: 13) AAC GGGUGAAGAUCCCCC (SEQ ID NO: 15) CCUGGGCCUG UGGCACCU (SEQ ID NO: 16) CCUGGGCCUG GGC CCUG (SEQ ID NO: 17) CCUGGGCCUG GGCA CUG (SEQ ID NO: 19) ACAGAGAAG GGGACUUC (SEQ ID NO: 20) ACAGAGAAG GGG GCUUCG (SEQ ID NO: 21) ACAGAGAAG GGGA UUCG (SEQ ID NO: 29)] GGA ACUCG CAGGUACCGC [0071] In a preferred embodiment of the invention, all of the above-listed siRNAs are provided, for in vivo use, with the short dAdT sequence protruding to the 3′ end. [0072] Method for the Preparation of the siRNAs [0073] Nucleotide synthesis methods for the preparation of short RNA sequences are known to a person skilled in the art and described in the state of the prior art. The siRNAs of the invention were produced by chemical synthesis, and are represented by duplexes of small oligonucleotides. These are comprised of 19 ribonucleotides with 2 deoxyribonucleotide “overhangs” at the 3′ end. Post-synthesis, the siRNAs were subjected to the following purification processes: Salts removal by ethanol precipitation, or using C 18-type chromatography columns Removal of 2′-ACE group present in RNA bases Pairing with the antisense sequence (synthesized in a separate reaction). Purification 1: the siRNA duplex is purified by acrylamide gel electrophoresis Purification 2: the siRNA duplex, obtained by the above-described step, is further purified with ion-exchange liquid chromatography (HPLC) Purification 3: the siRNA duplex, obtained by the above-described step, is subjected to counterionic exchange (Na + ), desalted, sterilized by filtration and tested for presence of endotoxins. [0080] Compositions and Dosages [0081] The siRNAs of the invention, their chemical derivatives and/or precursors can be administered systemically or locally. [0082] Tests conducted in vitro on cell cultures and in vivo on animal model demonstrated that the siRNAs of the invention are effectively internalized in the cell with no need of any transfection agent, rather by mere incubation with the cell in solution. In fact, the incubation, under standard cell culture conditions, of osteoclasts differentiated from peripheral blood mononuclear cells of healthy donors with siRNA of the invention or derivatives thereof highlighted siRNA incorporation into the cell and its preservation up to +7 days after treatment ( FIG. 10 ). [0083] Therefore, pharmaceutical compositions suitable to the administration of the siRNAs of the invention or of their chemical derivatives are compositions containing a pharmaceutically effective amount of siRNA, its derivative or its precursor, in a suitable, essentially liquid excipient. Such compositions are in the form of solutions, suspensions or emulsions. Any pharmaceutical excipient suitable for such applications can therefore be used. Suitable excipients are physiological solutions for parenteral use, hydroalcoholic solutions, glycol solutions, water/oil or oil/water emulsions, liposome or exosome emulsions/suspensions, oily solutions, micellar suspensions, vesicles, or complexes with PEI (polyethyleneimine) or complexes with atelocollagen, all containing the usual pharmaceutical additives, diluents, stabilizers and pH adjusters to physiological values. [0084] Administration of the siRNA of the invention, derivatives or precursors thereof, can occur parenterally, e.g., the intravenous, intraperitoneal, intramuscular, intradermal, subcutaneous, intraosseus, intracartilagineous, intraarticular administration. Alternatively, the administration can be carried out orally, through pills, tablets, formulations for buccal or sublingual dissolution, capsules, soft capsules, films, powders, granulate; rectally or vaginally, through suppositories or ovules; by inhalation, e.g. intrabronchial. [0085] Local administration can occur through any formulation suitable for local application, e.g. through topical application or direct application on or in the tissues to be treated, or again by local administration of a siRNA precursor and in situ production of the siRNA of the invention. Compositions based on exosomes, liposomes, vesicles, micelles containing the siRNA or their precursors are useful to attain both a systemic and a local effect. [0086] To obtain a local effect, the siRNAs of the invention or their derivatives or precursors can be administered through viral or nonviral vectors, or through the DNA encoding the siRNAs, or as isolated (naked) RNA (Pelled et al., 2010 Tissue Engineering: Part B, Volume 16, No. 1, 13-20) or through three-dimensional biocompatible matrices or implants, based, e.g., on fibrinogen and thrombin polymers and located in the application point. [0087] In a specific embodiment, the siRNAs or their derivatives or precursors are bound or associated or complexed to usual reagents ensuring an effective in vivo distribution of the siRNA, for instance polyethyleneimine (PEI) or derivatives thereof, such as the polyethyleneimine-polyethylene glycol-N-acetylgalactosamine (PEI-PEG-GAL) complex, or the polyethyleneimine-polyethylene glycol-tri-N-acetyl galactosamine (PEI-PEG-triGAL) complex. In a specific embodiment of the invention, the siRNAs are bound to the jetPEI® product, which is a linear derivative of polyethyeneimine provided by PolyPlus Transfection. [0088] Alternatively, the siRNAs of the invention can be locally administered in the form of their shRNA precursor within the scope of a gene therapy. For instance, a shRNA, or the DNA encoding a shRNA, can be transferred into a mammalian cell, by using, e.g., a suitable plasmid or an adenoviral vector as described by Egermann et al., Human Gene Ther . May 2006; 17 (5):507-17. The shRNAs expressed and processed by the cell itself produce the corresponding siRNAs able to silence the target gene. [0089] In an in vivo form of administration alternative to the vectors, the siRNAs can be transferred into a cell through electroporation, ultrasoundporation, cationic liposome-mediated transfection, microinjection, electropulsation. [0090] In another alternative form of local administration, the siRNAs of the invention, their derivatives or precursors, can be bound, adsorbed, immobilized even through covalent bonding to a matrix able to release the genetic material (gene-activated matrix (GAM)) as described by Luginbuehl et al., 2004, Eur J Pharm Biopharm 58:197-208, and then implanted in the zone of interest as described by Fang et al., 1996 (Proc Natl Acad Sci USA 93, 5753). [0091] Transfection agents, though not necessary, can however be used to improve siRNA internalization into osteoclasts. Transfection agents suitable for the present invention are: lipofectamine, nucleofection by Amaxa Nucleofector® (Lonza, Cologne, Germany) method using a specific kit (Cat# VPA-1007, Lonza). [0092] Posology [0093] Moreover, in vitro and in vivo tests conducted within the scope of the present invention demonstrated that the siRNAs internalized in the cell, i.e. in the osteoclasts, preserve their integrity and therefore their functionality over a period of several days. [0094] Hence, the treatment regimen with siRNAs of the invention provides administrations from once a day to once a week, e.g. 1, 2, 3, 4, 5, 6 or 7 administrations/week. Alternatively, the treatment can be carried out with a daily administration, or every 2, 3, 4, 5, 6, 7 days. [0095] The duration of the treatment depends on the severity of the disease and ranges from a treatment of some weeks to a chronic treatment. [0096] The tests carried out by the present inventors demonstrated that the siRNAs of the invention are effective in restoring osteoclast functionality in a broad spectrum of dosages, of from about 1 ng/kg of body weight to about 100 mg/kg of body weight of the subject to be treated, or subject in which symptoms of osteopetrosis progression have appeared. In a preferred embodiment, the dosages will be from about 1 μg/Kg to 20 mg/Kg of body weight, preferably from about 1 mg/Kg to about 10 mg/Kg. [0097] To be able to assay the in vivo efficacy of the siRNAs as potential medicaments, the following methodology was adopted: 1. It was verified that the siRNA for clcn7 normal mRNA were effective in reducing normal gene expression in WT mice. 2. It was verified that the siRNA against clcn7 mutated mRNA were not altering normal gene expression. 3. It was verified that the siRNA against clcn7 mutated mRNA were effective in reducing the mutated mRNA and in ameliorating the phenotype of ADO2 mice. [0101] The procedures for such verifications are described in the experimental examples. [0102] Combined Therapy [0103] The siRNAs of the invention can be used in association with other active principles. By the term “in association” it is meant both a co-therapy or combined therapy, and a co-formulation in a single pharmaceutical form, or in a single commercial package, e.g. a kit or a blister of two or more active principles. [0104] Active principles combinable with the siRNAs are for instance agents able to increase bone tissue anabolism: e.g. teriparatide, blosozumab, romosozumab, or even bone growth factors or nucleic acids encoding them, e.g., proteins of the BMP family, such as BMP-2 and/or BMP-7, or RNAs, like e.g. RNAs antagonizing the MIR-31, or transfection agents such as, e.g., lipofectamine, nucleofection by Amaxa Nucleofector® method (Lonza, Cologne, Germany) using a specific kit (Cat# VPA-1007, Lonza). EXPERIMENTAL SECTION Example 1: Generation of Vectors Carrying the Constructs of CLCN7 Gene Conjugated with the Sequence for EGFP (Enhanced Green Fluorescent Protein) [0105] To be able to perform the experiments in vitro, expression vectors were generated carrying the WT construct of the CLCN7 gene, conjugated with the EGFP sequence to allow visualization of the fusion protein by fluorescence analysis and quantification of the transcriptional expression of the gene by real-time RT-PCR for EGFP (WT CLCN7/pEGFP-C1) ( FIG. 1 ). For that purpose, the full sequence of human CLCN7 cDNA (http://www.ncbi.nlm.nih.gov/nuccore/NM_001114331.2) (rzpd IRAUp969B0859D6) was cloned in the pEGFP-C1 vector by restriction enzymes HindIII and XhoI. The full human sequence of the CLCN7 gene was amplified by iProof™ High-Fidelity DNA Polymerase kit (BIO-RAD 172-5301) using primers with the 5′ end provided with the sequences of restriction enzymes HindIII and XhoI. Then, a double digestion of the empty pEGFP-C1 vector and of the PCR product of CLCN7 was performed for 3 hours at 37° C., using the restriction enzymes HindIII and XhoI. The digested vector and PCR product were purified by QIAquick PCR Purification Kit (Qiagen 28104). Then, dephosphorylation of the digested vector was performed, for 1 hour at 37° C., using Shrimp Alkaline Phosphatase (SAP). 300 ng of the digested CLCN7 PCR product and 100 ng of the dephosphorylated pEGFP-C1 vector were ligated by T4 DNA ligase, overnight at 4° C. Ligation was then used to transform XIBlue1 cells. Subsequently, vectors with mutated construct p.R767W-, p.G215R, -p.A788D- and p.R286W-CLCN7/pEGFP-C1 were obtained by QuikChange II XL Site-Directed Mutagenesis Kit (Cat.#6200521, Stratagene), using primers containing the desired mutation. Occurred mutagenesis was then checked by direct DNA sequencing ( FIG. 2 ). Example 2: Transfections [0106] The vectors were used to transfect human HEK293 cells by standard transfection technique with lipofectamine, then expression of the corresponding mRNA was quantified by real-time RT-PCR. Transfections of the WT construct and of the mutated constructs induced similar levels of transcriptional expression ( FIG. 3 ). [0107] Then, expression of WT- and mutated proteins in transfected HEK293 cells was assessed by confocal microscopy for the detection of fusion proteins EGFP. Correct localization of the fluorescent protein was demonstrated by co-localization of Golgi apparatus markers (gamma-adaptin), early endosomes (EEA-1), acid vesicles (Lysotracker Red) and lysosomes (lamp-1) ( FIG. 4 ). [0108] To demonstrate the ability of the mutated constructs of the Inventors to be expressed in the osteoclast line and inhibit bone resorption, the murine line of RAW264.7 osteoclast precursors was transfected with empty vector, or with a vector carrying the mutated constructs. mRNA and protein expressions were checked by real-time RT-PCR and confocal microscopy, respectively, then the cells were plated onto bovine bone slices and differentiated into mature osteoclasts by treatment with 100 ng/ml RANKL for 4 days. Bone resorption was quantified by count of the resorption lacunae (pits) dug out by osteoclasts (pit index assay). The results demonstrated a bone resorption reduction of about 70% in RAW264.7 cells transfected with the mutated constructs, compared to the same cells transfected with the empty vector ( FIG. 5 ). This percentage of bone resorption is very similar to that observed in osteoclasts differentiated from the peripheral blood of ADO2 patients, compared to osteoclasts from healthy donors (6). [0109] These experiments yielded good evidence that the generated vectors might represent valid tools for evaluating the Inventors' strategy of in vitro silencing of the mutated gene. However, to verify that the method would work also in primary human osteoclasts, a nucleofection process by AMAXA nucleofector was set up. The procedure was successful, enabling a good hyperexpression of the Inventor's EGFP fusion protein in osteoclasts differentiated from peripheral blood mononuclear cells of healthy donors ( FIG. 6 ). [0110] Assessment of bone resorption in human osteoclasts transfected with the human constructs demonstrated once more a 70% reduction compared to cells transfected with the empty vector (not shown). Example 3: In Vitro Treatments with siRNA [0111] After having set up the investigation methods, siRNAs against the abovementioned mutations of the CLCN7 gene (Table 2) were designed and assigned for synthesis to Dharmacon Company. Moreover, the commercial pool of siRNA against WT CLCN7 gene and scrambled (mixed nucleotide sequence) control siRNA was purchased. [0112] The setting up of the procedure was performed with siRNA against p.R767W mutation. It was then extended to the other mutations. A siRNA specific for the trascriptor carrying the p.R767W mutation (R286W siRNA 1) was assessed for efficacy and specificity in HEK293 cells transfected with WT- or p.R767W-CICN7/EGFP vectors. The results showed a ≦60% reduction of mutated mRNA expression. However, this positive result was invalidated by a similar reduction of normal mRNA in cells transfected with the WT CLCN7/EGFP vector ( FIG. 7 ). [0113] Although this siRNA showed no specificity for the mutated gene, its good efficacy in reducing the mRNA of the mutated transcript was encouraging. Therefore, its effect on bone resorption was assessed by using RAW264.7 cells transfected with the R767W- or WT-CLCN7/EGFP construct. Under these experimental conditions, it was demonstrated that the siRNA for the mutated mRNA showed a partial ability to reactivate bone resorption compared to control scrambled siRNA. In this experiment, a modest inhibition of bone resorption was observed in cells transfected with WT-CLCN7/EGFP vector and subjected to treatment with siRNA against p.R767W mutation ( FIG. 8 ). This reduction, lower than what observed in cells transfected with mutated vector (compare second and third bar from the left) might be due to abundance of CLCN7 mRNA expression in RAW264.7 cells, due to the presence both of endogenous mRNA and of mRNA produced by the transfected construct. [0114] At this point, a strategy for increasing the specificity toward the mutated mRNA was adopted. This was obtained by inserting a non-complementary (mismatch) nucleotide in various positions downstream of the mutated nucleotide (11). With this strategy, three novel siRNA for p.R767W mutation were designed (Table 2). All three siRNAs showed greater specificity for the mutated mRNA compared to the WT mRNA. With one of them (termed p.R767W siRNA 2C), a 90% efficacy of mutated mRNA reduction was obtained, without any reduction of WT mRNA ( FIG. 9 ). [0115] Then, human osteoclasts were treated with siRNA for the normal CLCN7 gene using siRNAs conjugated with Cy3 fluorophore (Cy3-WT siRNA). The aim was to set up the strategy of siRNA internalization into primary cells, using confocal microscopy for its checking. Under these conditions, it was observed that osteoclasts internalize the siRNAs with no need of any transfection agent. In fact, incubation, under standard culture conditions, of osteoclasts differentiated from peripheral blood mononuclear cells of healthy donors with 300 nM Cy3-WT siRNA highlighted siRNA incorporation into the cell and its preservation up to +7 days from treatment ( FIG. 10A ). [0116] To demonstrate the ability of R767W siRNA C to reduce CLCN7 mutated p.R767W expression and decrease its detrimental effect on bone resorption, human osteoclasts were transfected with the p.R767W CLCN7/pEGFP vector and treated with 500 nM p.R767W siRNA 2C. Under these circumstances, there were highlighted a reduced transcriptional expression of the EGFP fluorescent protein sequence ( FIG. 10B ) and the restoration of osteoclasts' ability to resorb bone ( FIG. 100 ) in cells treated with p.R767W siRNA 2C, compared to osteoclasts treated with scrambled siRNA. [0117] Then, it was evaluated whether the treatment with p.R767W siRNA 2C influenced the transcriptional expression of the normal CLCN7 transcript and the ability to resorb bone in human osteoclasts from a healthy donor. The results showed a good effectiveness of the siRNA directed against normal CLCN7 (used as positive control) in decreasing this mRNA, whereas no effect thereon by p.R767W siRNA 2C was observed ( FIG. 11 ). In accordance with this result, p.R767W siRNA 2C did not modify the bone resorption of osteoclasts from healthy donors ( FIG. 12 ). [0118] Overall, these results show that the strategy of the Inventors was successful and allowed them to design highly specific siRNAs against the p.R767W mutation of the CLCN7 gene, which had no effect on the normal mRNA of human osteoclasts. [0119] Following the same strategy, siRNAs against the other aforedescribed three mutations, p.G215R, p.A788D and p.R286W (Table 2) were designed. From a detailed analysis of the results, it emerged that for the p.A788D mutation siRNAs meeting the efficacy and specificity criteria required for their use in the therapy of CLCN7-dependent ADO2 have not been identified yet. As to mutations p.R286W and p.G215R, siRNAs were instead identified for which the efficacy and specificity criteria were met (Table 2). Since for the p.G215R siRNA mutation a murine model of disease is presently available, the features of the effective siRNA identified in Table 2 with the abbreviation p.G215R 2M were studied. It proved highly active in reducing the expression of mutated mRNA in HEK293 cells and in primary human osteoclasts ( FIG. 13 , FIG. 28 ). [0120] Moreover, by using murine osteoclasts from the animal ADO2 model generated in the laboratory of the Inventors (10), carrying the murine homologue (p.G213R) of the human mutation (p.G215R), it was demonstrated that the treatment with p.G213R siRNA (Table 3) was able to increase bone resorption ( FIG. 14 ). Example 4: In Vivo Treatments [0121] To be able to assay siRNAs efficacy in vivo, as medicaments, the following procedure was followed: [0122] i) Verifying that the siRNA for clcn7 normal mRNA be effective in reducing normal gene expression in WT mice. [0123] ii) Verifying that the siRNA against clcn7 mutated mRNA does not alter normal gene expression. [0124] iii) Verifying that the siRNA against clcn7 mutated mRNA be effective in reducing the mutated mRNA and in ameliorating ADO2 mice phenotype. [0125] i) Verifying that the siRNA for Clcn7 Normal mRNA be Effective in Reducing Normal Gene Expression in WT Mice. [0126] This group of experiments was carried out with a pool of siRNAs against the normal clcn7 gene, available on the market by Dharmacon, whose nucleotide sequence is unknown. 1-month old CD1 mice (n=4) were treated with 2 mg/Kg clcn7-siRNA inoculated by TransIT-QR (Quick Recovery) Hydrodynamic Delivery kit. The TransIT solution is specifically studied for safe and effective administration of nucleic acids, using the hydrodynamic injection procedure in the tail vein. After 24-48 hours, animals were sacrificed and subjected to anatomical dissection to collect heart, spleen, liver, kidneys, brain, lungs and tibias ( FIG. 15 ). This latter result probably depends on clcn7-siRNA inability to cross the blood-brain barrier. [0127] Then, in normal CD1 mice it was also ascertained which were the best clcn7-siRNA administration pathway, in order to verify the feasibility of repeated treatments. Injection by TransIT-QR Hydrodynamic Delivery kit proved effective, but it can be performed only once and on adult mice. Intraperitoneal administration proved easy to perform from the animals' first days of life. Venous infusion was also efficient, but it was possible to perform it only on adult mice and for a very limited number of times. Therefore, intraperitoneal injection was chosen for subsequent studies. Moreover, the best administration frequency was tested to be of 48 hours, by which the best dose-dependent response was demonstrated both in bone and in other organs. The maximum dose used in this series of experiments (0.5 mg/Kg of body weight), administered 3 times a week for 3 weeks, induced no modifications of structural bone parameters, measured by computerized microtomography (μCT). The Inventors explain this negative result with the notion that the gene is haplosufficient and that this treatment regimen has reduced the clcn7 mRNA only of 60%, leaving a 40% of mRNA probably sufficient for the carrying out of its functions. In any case, the treatment induced no sign of suffering, nor did it cause any evident distress to the animals. [0128] ii) Verifying that the siRNA Against Clcn7 Mutated mRNA does not Alter Normal Gene Expression and does not Induce Adverse Effects. [0129] To demonstrate this aspect, normal CD1 mice were treated with p.G213R-clcn7 siRNA (Table 3) (0.5 mg/Kg) and it was observed that there was no reduction of the normal transcript, unlike what found in the treatment with the clcn7-siRNA directed against the normal transcript. In all these experiments, the control scrambled siRNA never caused alterations of the expression of the clcn7 gene, neither normal, nor mutated. With these experiments the Inventors therefore demonstrated that the siRNAs for one of the mutations of the clcn7 gene are ineffective towards the normal transcript and do not induce adverse side effects. [0130] iii) Verifying that the siRNA Against Clcn7 Mutated mRNA be Effective in Reducing the Mutated mRNA and in Improving ADO2 Mice Phenotype [0131] To test the effectiveness of their treatment, an experiment was therefore carried out in the sole murine ADO2 model available (12). This model was created in C57BL/6 mice strain by knock-in technology, which allowed substitution of the normal hexon 7 of the clcn7 gene with a hexon 7 mutated by a G-A transition in position 14365 of the DNA, corresponding to the protein mutation p.G213R. [0132] Mice homozygous for this mutation are small, lack teeth eruption and die within 30 days from birth, even when fed a soft diet. They exhibit an extremely severe osteopetrotic phenotype, fibrous bone marrow and hippocampal and cerebellar cortex degeneration similar to what found in a clcn7 knock-out murine model and in human autosomal recessive osteopetrosis. [0133] Heterozygous mice are born at the normal Mendelian frequency, are vital and fertile and unaffected by alterations of size, body weight and teeth eruption. They instead exhibit the typical signs of a less severe osteopetrosis, without evident signs of neurodegeneration. Heterozygous adult (3 month-old) mice exhibit greater mineral density and greater bone mass, verified by μCT analysis of trabecular structural bone parameter of the tibias, femurs and vertebrae. This increase of bone mass is persistent, can also be found in old mice, and is similar in males and in females. Histological examination of heterozygous 3 month-old mice showed an increase of the expression of the osteoclast-specific enzyme TRAcP and an increase of osteoclast number/surface/bone surface. Despite this increase, bone resorption is reduced as indicated by serum levels of bone resorption marker CTX normalized for serum activity of osteoclastic enzyme TRAcP. On the contrary, all bone formation parameters [serum marker (osteocalcin), osteoblast surface/bone surface, bone formation rate, osteoid thickness, growth plate thickness] demonstrate that there is no osteoblast or chondrocyte involvement, nor do mice have a phenotype compatible with osteopetrorachitism. [0134] Bone marrow collected from mice shows an increase of the number of osteoclast precursors and a greater osteoclastogenesis in vitro in the presence of M-CSF and RANKL. Nevertheless, bone resorption is reduced compared to osteoclasts obtained from the bone marrow of normal mice. Heterozygous mice have normal hematological and serum parameters (pancreatic amylase, hepatic transaminases, potassium, calcium, phosphorus, muscle creatine kinase and glucose concentration), whereas parathyroid hormone levels are increased, in accordance with increased osteoclastogenesis. [0135] Having obtained a reliable murine ADO2 model, therein it was verified whether the therapy with siRNAs directed against mutation p.G215R were effective. First of all, optimal dose and administration time of p.G213R-clcn7 siRNA were established. To this end, primers able to amplify exclusively the mutated transcript ( FIGS. 16 and 17 ) were first designed. Then, p.G213R-clcn7 ADO2 mice were treated with 2 or 4 mg/Kg of body weight of p.G213R-clcn7 sticky siRNA/jetPEI by intraperitoneal injection, verifying the in-serum kinetics of the total RNA ( FIG. 18 ) and confirming the reduction of p.G213R-clcn7 mutated mRNA expression by real-time RT-PCR in mice treated with p.G213R-clcn7 sticky siRNA/jetPEI compared to mice treated with control (scrambled) siRNA ( FIG. 19 ). This reduction was also confirmed in tibia ( FIG. 20 ) and was not evident anymore after 96 hours from p.G213R-clcn7 sticky siRNA/jetPEI administration, a circumstance indicating the best treatment frequency to be of 48 hours. [0136] To verify whether this treatment might have an effect on bone resorption in vivo, 10 day-old p.G213R-clcn7 ADO2 mice were treated, 3 times a week for 2 weeks and 4 weeks, with 4 mg/Kg of body weight of control (scrambled) sticky siRNA/jetPEI, or p.G213R-clcn7-specific siRNA. The treatment was well-tolerated and did not induce histopathological damages to vital organs ( FIG. 21 ). Moreover, it improved renal, hepatic and muscle damage biomarkers ( FIG. 22 ). The therapeutic effect of the treatment was analyzed on bone resorption biomarker CTX, normalized for osteoclastic biomarker TRAcP. The results demonstrated a significant increase of CTX serum levels and of the CTX/TRAcP ratio in mice treated with p.G213R-clcn7 sticky siRNA/jetPEI, as evidence of occurred activation of osteoclastic bone resorption ( FIG. 23 ). Consistently, μCT analysis of the proximal end of the tibia showed a reduction of the trabecular bone volume/total tissue volume percentage and an improvement of structural trabecular variables already after 2 weeks of treatment ( FIG. 24 ). After 4 weeks, total restoration of structural parameters ( FIG. 25 ) and osteoclast functionality ( FIG. 26 ) were witnessed, without any undesired effect on osteoblastic parameters ( FIG. 27 ). Finally, an improvement also of biomechanic parameters was witnessed ( FIG. 28 ), indicative of restoration of a good bone tissue quality. These results indicate that the bone resorption increase induced by the treatment of the Inventors was effective in correcting the bone phenotype of ADO2 mice. [0137] To complete the therapeutic study, the ability of the identified siRNAs in reducing human mutated CLCN7 expression in osteoclasts from healthy donors transfected with the constructs carrying the mutations p.G215R, p.R286W and p. R767W, as well as the ability of the p.G215R-specific siRNA to improve bone resorption in osteoclasts obtained from a patient ( FIG. 29 ) were further confirmed. [0138] As to the siRNAs for the p. R767W mutation and for any other mutation to be analyzed in vivo, at present it cannot be suggested that murine models may be generated for each mutation, considering both the times and costs of implementation. Therefore, to test the in vivo efficacy of p.R767W siRNA 2C, which the Inventors had found to be active in vitro, an alternative strategy was adopted. Human breast cancer MDA-MB-231 cells were stably transfected with p.R767W-CLCN7/EGFP-C1 vector and treated in vitro with R767W siRNA 2C to verify its efficacy and specificity on the silencing of CLCN7 mutated mRNA ( FIG. 30A ). Then, the cells transfected with the p.R767W-CLCN7/EGFP-C1 vector were injected in the subcutaneous tissue of athymic (immunocompromised) Balb/c nu/nu mice, in which they formed macroscopically evident tumors. When the tumors reached the volume of 1 cm 3 , the mice were treated with a single intraperitoneal injection of 4 mg/Kg of body weight of R767W 2C/jetPEI siRNA. After 48 hours mice were sacrificed, tumors were excided and analyzed, by real-time RT-PCR, for the expression of the EGFP transcript conjugated to the mutated gene. Under these experimental conditions, a tumor-expressed p.R767W-CLCN7/EGFP transcriptional reduction of about 50% was obtained. To check the specificity of the siRNA contrived by the Inventors, the same treatment protocol was carried out also by using siRNA for the WT CLCN7 gene, finding no change of transcriptional expression of the EGFP conjugated with the mutated construct ( FIG. 30 ). REFERENCE [0000] 1. Deng Y, Wang C C, Choy K W, Du Q, Chen J, Wang Q, Chung T K, Tang T. Therapeutic potentials of gene silencing by RNA interference: Principles, challenges, and new strategies. Gene. 538:217-227, 2014. 2. J. Bollerslev. Osteopetrosis, A genetic and epidemiological study, Clin. Genet. 31:86-90, 1987. 3. O. D. Bénichou, J. D. Laredo, M. C. de Vernejoul, Type II autosomal dominant osteopetrosis (Albers-Schönberg disease): clinical and radiological manifestations in 42 patients. Bone 26 (2000) 87-93. 4. A. Frattini, A. Pangrazio, L. Susani, C. Sobacchi, M. Mirolo, M. Abinun, M. Andolina, A. Flanagan, E. M. Horwitz, E. Mihci, L. D. Notarangelo, U. Ramenghi, A. Teti, J. Van Hove, D. Vujic, T. Young, A. Albertini, P. J. Orchard, P. Vezzoni, A. Villa, Chloride channel CICN7 mutations are responsible for severe recessive, dominant, and intermediate osteopetrosis. J. Bone Miner. Res. 18:1740-1747, 2003. 5. B. Peruzzi, A. Teti. The physiology and pathophysiology of the osteoclast. Clin Rew Bone Miner Metab. 10:71-97, 2012. 6. A. Del Fattore, B. Peruzzi, N. Rucci, I. Recchia, A. Cappariello, M. Longo, D. Fortunati, P. Ballanti, M. lacobini, M. Luciani, R. Devito, R. Pinto, M. Caniglia, E. Lanino, C. Messina, S. Cesaro, C. Letizia, G. Bianchini, H. Fryssira, P. Grabowski, N. Shaw, N. Bishop, D. Hughes, R. P. Kapur, H. K. Datta , A. Taranta, R. Fornari, S. Migliaccio, A. Teti. Clinical, genetic, and cellular analysis of 49 osteopetrotic patients: implications for diagnosis and treatment, J. Med. Genet. 43:315-325, 2006. 7. K. Chu, D. L. Koller, R. Snyder, T. Fishburn, D. Lai, S. G. Waguespack, T. Foroud, M. J. Econs. Analysis of variation in expression of autosomal dominant osteopetrosis type 2: searching for modifier genes, Bone 37:655-661, 2005. 8. E. Cleiren, O. Bénichou, E. Van Hul, J. Gram, J. Bollerslev, F. R. Singer, K. Beaverson, A. Aledo, M. P. Whyte, T. Yoneyama, M. C. deVernejoul, W. Van Hul, Albers-Schönberg disease (autosomal dominant osteopetrosis, type II) results from mutations in the CICN7 chloride channel gene, Hum. Mol. Genet. 10:2861-2867, 2001. 9. U. Kornak, D. Kasper, M. R. Bösl, E. Kaiser, M. Schweizer, A. Schulz, W. Friedrich, G. Delling, T. J. Jentsch, Loss of the CIC-7 chloride channel leads to osteopetrosis in mice and man. Cell: 104:205-215, 2001. 10. A. Del Fattore, M. Capannolo, N. Rucci, A. Teti. New experimental therapeutic approach by siRNA for autosomal dominant osteopetrosis (ADO). Bone 46:142, 2010 (abstract). 11. Y. Ohnishi, Y. Tamura, M. Yoshida, K. Tokunaga, H. Hohjoh H. Enhancement of allele discrimination by introduction of nucleotide mismatches into siRNA in allele-specific gene silencing by RNAi. PLoS ONE. 3:e2248, 2008. 12. I. Alam, A. K. Gray, K. Chu, S. Ichikawa, K. S. Mohammad, M. Capannolo, M. Capulli, A. Maurizi, M. Muraca, A. Teti, M. J. Econs, A. Del Fattore. Generation of the first autosomal dominant osteopetrosis type II (ADO2) disease models. Bone. 59:66-75, 2014.

The present invention lies in the field of molecules known as “small interfering RNA” with therapeutic applications. siRNAs have the ability to reduce gene expression in an extremely specific way (1). These are small sequences of double-strand RNA, normally used in laboratory to modify cell function, which revolutionized cell biology by allowing previously precluded molecular manipulations.

CROSS-REFERENCE TO RELATED APPLICATIONS [0001] This application claims the benefit of priority of Korean Patent Application No. 10-2011-0140861 filed on Dec. 23, 2011, Korean Patent Application No. 10-2012-0003617 filed on Jan. 11, 2012, and Korean Patent Application No. 10-2012-0147996 filed on Dec. 18, 2012, all of which is incorporated by reference in its entirety herein. BACKGROUND OF THE INVENTION [0002] 1. Field of the Invention [0003] The present invention relates to an image processing method and apparatus and, more particularly, to an inter-frame prediction method and an apparatus using the method. [0004] 2. Related Art [0005] A demand for images having high resolution and high quality, such as a High Definition (HD) image and an Ultra High Definition (UHD) image, is recently increasing in a variety of application fields. As the resolution and quality of image data become higher, the amount of the image data becomes relatively greater than that of the existing image data. For this reason, if the image data is transmitted using media, such as the existing wired/wireless broadband lines, or the image data is stored by using the existing storage medium, a transmission cost and a storage cost are increased. Image compression techniques with high efficiency can be used to solve the problems occurring as the resolution and quality of image data becomes higher. [0006] Image compression techniques include a variety of techniques, such as an inter-frame prediction technique for predicting a pixel value included in a current picture from a picture anterior or posterior to the current picture, an intra-frame prediction technique for predicting a pixel value included in a current picture by using information on a pixel within the current picture, and an entropy coding technique for allocating a short symbol to a value having high frequency of appearance and allocating a long symbol to a value having low frequency of appearance. Image data can be effectively compressed, transmitted, or stored by using the image compression techniques. SUMMARY OF THE INVENTION [0007] An object of the present invention is to provide a method of setting the reference picture index of a temporal merging candidate. [0008] Another object of the present invention is to provide an apparatus for performing a method of setting the reference picture index of a temporal merging candidate. [0009] In accordance with an aspect of the present invention, an inter-frame prediction method using a temporal merging candidate may include the steps of determining the reference picture index of the temporal merging candidate for a current block and deriving the temporal merging candidate block of the current block and deriving the temporal merging candidate from the temporal merging candidate block, wherein the reference picture index of the temporal merging candidate can be derived irrespective of whether other blocks except the current block have been decoded or not. The temporal merging candidate may be derived in a unit of a coding block including the current block or in a unit of the current block depending on whether the current block will use a single merging candidate list or not. The inter-frame prediction method may further include the step of determining whether or not the current block is a block using the single merging candidate list, wherein the single merging candidate list may derive and generate at least one of the spatial merging candidate and the temporal merging candidate of a prediction block based on a coding block including the prediction block. The step of determining whether or not the current block is a block using the single merging candidate list may include the steps of decoding information on the size of the current block and determining whether or not the information on the size of the current block satisfies conditions of the size of a block that the single merging candidate list is derived. The reference picture index of the temporal merging candidate may be set to a fixed value. The temporal merging candidate may include a temporal motion vector calculated by comparing a difference between the reference picture index of a temporal merging candidate block (i.e., a colocated block) and the index of a picture (i.e., a colocated picture) including the colocated block with a difference between the reference picture index of the temporal merging candidate having the index of the fixed value and the index of the picture including the current block. The reference picture index of the temporal merging candidate may be set to 0. [0010] In accordance with another aspect of the present invention, a decoder for performing an inter-frame prediction method using a temporal merging candidate includes a merging candidate deriving unit configured to determine the reference picture index of the temporal merging candidate for a current block, derive the temporal merging candidate block of the current block, and derive a temporal merging candidate from the temporal merging candidate block, wherein the reference picture index of the temporal merging candidate may be derived irrespective of whether other blocks except the current block have been decoded or not. The temporal merging candidate may be derived in a unit of a coding block including the current block or in a unit of the current block depending on whether the current block will use a single merging candidate list or not. The merging candidate deriving unit may be configured to determine whether or not the current block is a block using the single merging candidate list, and the single merging candidate list may derive and generate at least one of the spatial merging candidate and the temporal merging candidate of a prediction block based on a coding block including the prediction block. The merging candidate deriving unit may be configured to decode information on the size of the current block and determine whether or not the information on the size of the current block satisfies conditions of the size of a block that the single merging candidate list is derived, in order to determine whether or not the current block is a block using the single merging candidate list. The reference picture index of the temporal merging candidate may be set to a fixed value. The temporal merging candidate may include a temporal motion vector calculated by comparing a difference between the reference picture index of a temporal merging candidate block (a colocated block) and the index of a picture (a colocated picture) including the colocated block with a difference between the reference picture index of the temporal merging candidate having the index of the fixed value and the index of the picture including the current block. The reference picture index of the temporal merging candidate may be set to 0. [0011] As described above, in accordance with the method and apparatus for setting the reference picture index of a temporal merging candidate according to embodiments of the present invention, inter-frame prediction using a temporal merging candidate can be performed on a plurality of prediction blocks in parallel by using a temporal merging candidate set to a specific value or using the reference picture index of a spatial merging candidate at a predetermined location as the reference picture index of a temporal merging candidate. Accordingly, an image processing speed can be increased, and the complexity of image processing can be reduced. BRIEF DESCRIPTION OF THE DRAWINGS [0012] FIG. 1 is a block diagram showing the construction of an image coder in accordance with an embodiment of the present invention. [0013] FIG. 2 is a block diagram showing the construction of an image decoder in accordance with another embodiment of the present invention. [0014] FIG. 3 is a conceptual diagram illustrating an inter-frame prediction method using merge mode in accordance with an embodiment of the present invention. [0015] FIG. 4 is a conceptual diagram illustrating inter-frame prediction using a temporal merging candidate and the reference picture index of the temporal merging candidate in accordance with an embodiment of the present invention. [0016] FIG. 5 is a conceptual diagram illustrating a case where one coding block is partitioned into two prediction blocks. [0017] FIG. 6 is a conceptual diagram illustrating a method of setting the reference picture index of a temporal merging candidate in accordance with an embodiment of the present invention. [0018] FIG. 7 is a conceptual diagram illustrating a method of deriving the reference picture indices of temporal merging candidates in accordance with an embodiment of the present invention. [0019] FIG. 8 is a conceptual diagram illustrating a method of deriving the reference picture indices of temporal merging candidates in accordance with an embodiment of the present invention. [0020] FIG. 9 is a conceptual diagram illustrating a method of deriving the reference picture indices of temporal merging candidates in accordance with an embodiment of the present invention. [0021] FIG. 10 is a conceptual diagram illustrating a method of deriving the reference picture indices of temporal merging candidates in accordance with an embodiment of the present invention. [0022] FIG. 11 is a conceptual diagram illustrating a method of deriving the reference picture index of a temporal merging candidate in accordance with an embodiment of the present invention. [0023] FIG. 12 is a conceptual diagram illustrating a method of deriving the reference picture index of a temporal merging candidate in accordance with an embodiment of the present invention. [0024] FIG. 13 is a conceptual diagram illustrating a method of deriving the reference picture index of a temporal merging candidate in accordance with an embodiment of the present invention. [0025] FIG. 14 is a flowchart illustrating a method of including a temporal merging candidate in a merging candidate list in accordance with an embodiment of the present invention. [0026] FIG. 15 is a conceptual diagram illustrating a method of generating a single merging candidate list by sharing all spatial merging candidates and temporal merging candidates in a plurality of prediction blocks in accordance with an embodiment of the present invention. [0027] FIG. 16 is a conceptual diagram illustrating a method of generating a single candidate list in accordance with an embodiment of the present invention. DESCRIPTION OF EXEMPLARY EMBODIMENTS [0028] Hereinafter, exemplary embodiments are described in detail with reference to the accompanying drawings. In describing the embodiments of the present invention, a detailed description of the known functions and constructions will be omitted if it is deemed to make the gist of the present invention unnecessarily vague. [0029] When it is said that one element is “connected” or “coupled” to the other element, the one element may be directly connected or coupled to the other element, but it should be understood that a third element may exist between the two elements. Furthermore, in the present invention, the contents describing that a specific element is “included (or comprised)” does not mean that elements other than the specific element are excluded, but means that additional elements may be included in the implementation of the present invention or in the scope of technical spirit of the present invention. [0030] Terms, such as the first and the second, may be used to describe various elements, but the elements should not be restricted by the terms. The terms are used to only distinguish one element and the other element from each other. For example, a first element may be named a second element without departing from the scope of the present invention. Likewise, a second element may also be named a first element. [0031] Furthermore, elements described in the embodiments of the present invention are independently shown in order to indicate different and characteristic functions, and it does not mean that each of the elements consists of separate hardware or a piece of software unit. That is, the elements are arranged, for convenience of description, and at least two of the elements may be combined to form one element or one element may be divided into a plurality of elements and the plurality of elements may perform functions. An embodiment in which the elements are combined or each of the elements is divided is included in the scope of the present invention without departing from the essence of the present invention. [0032] Furthermore, in the present invention, some elements may not be essential elements for performing essential functions, but may be optional elements for improving only performance. The present invention may be implemented using only the essential elements for implementing the essence of the present invention other than the elements used to improve only performance, and a structure including only the essential elements other than the optional elements used to improve only performance are included in the scope of the present invention. [0033] FIG. 1 is a block diagram showing the construction of an image coder in accordance with an embodiment of the present invention. [0034] Referring to FIG. 1 , the image coder 100 includes a motion prediction unit 111 , a motion compensation unit 112 , an intra-prediction unit 120 , a switch 115 , a subtractor 125 , a transform unit 130 , a quantization unit 140 , an entropy coding unit 150 , an inverse quantization unit 160 , an inverse transform unit 170 , an adder 175 , a filtering unit 180 , and a reference picture buffer 190 . [0035] The image coder 100 can perform coding an input picture in intra mode or inter mode and output a bit stream. The switch 115 can switch to intra mode in the case of intra mode and can switch to inter mode in the case of inter mode. The image coder 100 can derive a prediction block for the input block of the input picture and then code the residual of the input block and the prediction block. [0036] Intra mode can be defined and used as a term ‘intra-frame prediction mode’, inter mode can be defined and used as a term ‘inter-frame prediction mode’, the intra-prediction unit 120 can be defined and used as a term ‘intra-frame prediction unit’, and the motion prediction unit 111 and the motion compensation unit 112 can be defined and used as a term ‘inter-frame prediction unit’. [0037] An inter-frame prediction method in accordance with an embodiment of the present invention discloses a method of determining the reference picture index of a temporal merging candidate. The intra-prediction unit 120 can include a merging candidate deriving unit for deriving the spatial merging candidate and temporal merging candidate blocks of a current block and deriving a spatial merging symbol from the spatial merging candidate block and a temporal merging candidate from the temporal merging candidate block. A method of deriving the merging candidates will be described in detail later. [0038] In the case of intra mode, the intra-prediction unit 120 can derive the prediction block by performing spatial prediction by using the pixel value of an already coded block near a current block. [0039] In the case of inter mode, the motion prediction unit 111 can obtain a motion vector by searching a reference picture, stored in the reference picture buffer 190 , for a region that is most well matched with the input block in a motion prediction process. The motion compensation unit 112 can derive the prediction block by performing motion compensation using the motion vector. [0040] The subtractor 125 can derive a residual block by way of the residual of the input block and the derived prediction block. The transform unit 130 can output a transform coefficient by performing transform on the residual block. Here, the transform coefficient can mean a coefficient value derived by performing transform on the residual block and/or a residual signal. In the following specification, a quantized transform coefficient level derived by applying quantization to a transform coefficient can also be called a transform coefficient. [0041] The quantization unit 140 can quantize the input transform coefficient according to a quantization parameter and output a quantized transform coefficient level. [0042] The entropy coding unit 150 can perform entropy coding based on values calculated by the quantization unit 140 or a coding parameter value derived in a coding process and output a bit stream based on a result of the entropy coding. [0043] If entropy coding is applied, the size of a bit stream for each of target coding symbols can be reduced because the symbols are represented by allocating a small number of bits to a symbol having a high probability of occurrence and a large number of bits to a symbol having a low probability of occurrence. Accordingly, the compression performance of image coding can be increased by way of the entropy coding. The entropy coding unit 150 can use a coding method, such as exponential golomb, Context-Adaptive Variable Length Coding (CAVLC), or Context-Adaptive Binary Arithmetic Coding (CABAC), for the entropy coding. [0044] In the image coder according to the embodiment of FIG. 1 , a currently coded image needs to be decoded and stored in order to be used as a reference picture because inter-prediction coding, that is, inter-frame prediction coding, is performed. Accordingly, the quantized coefficient is inversely quantized by the inverse quantization unit 160 and then inversely transformed by the inverse transform unit 170 . The inversely quantized and inversely transformed coefficient is added to the prediction block by way of the adder 175 , and thus a reconstructed block is derived. [0045] The reconstructed block experiences the filtering unit 180 . The filtering unit 180 can apply one or more of a deblocking filter, a Sample Adaptive Offset (SAO), and an Adaptive Loop Filter (ALF) to the reconstructed block or a reconstructed picture. The reconstructed block that has experienced the filtering unit 180 can be stored in the reference picture buffer 190 . [0046] FIG. 2 is a block diagram showing the construction of an image decoder in accordance with another embodiment of the present invention. [0047] Referring to FIG. 2 , the image decoder 200 includes an entropy decoding unit 210 , an inverse quantization unit 220 , an inverse transform unit 230 , an intra-prediction unit 240 , a motion compensation unit 250 , an adder 255 , a filtering unit 260 , and a reference picture buffer 270 . [0048] The image decoder 200 can receive a bit stream from a coder, perform decoding on the bit stream in intra mode or inter mode, and output a reconfigured image, that is, a reconstructed picture. A switch can switch to intra mode in the case of intra mode and can switch to inter mode in the case of inter mode. The image decoder 200 can obtain a reconstructed residual block from the input bit stream, derive a prediction block from the reconstructed residual block, and derive a block reconstructed by adding the reconstructed residual block and the prediction block together, that is, the reconstructed block. [0049] The entropy decoding unit 210 can derive symbols, including a symbol having a quantized coefficient form, by performing entropy decoding on the input bit stream according to a probability distribution. The entropy decoding method is similar to the aforementioned entropy coding method. [0050] If the entropy decoding method is applied, the size of a bit stream for each of symbols can be reduced because the symbols are represented by allocating a small number of bits to a symbol having a high probability of occurrence and a large number of bits to a symbol having a low probability of occurrence. Accordingly, the compression performance of image decoding can be increased by the entropy decoding method. [0051] The quantized coefficient is inversely quantized by the inverse quantization unit 220 and then inversely transformed by the inverse transform unit 230 . As a result of the inverse quantization and the inverse transform on the quantized coefficient, the reconstructed residual block can be derived. [0052] In the case of intra mode, the intra-prediction unit 240 can derive a prediction block by performing spatial prediction using the pixel value of an already decoded block near a current block. In the case of inter mode, the motion compensation unit 250 can derive the prediction block by performing motion compensation using a motion vector and a reference picture stored in the reference picture buffer 270 . [0053] An inter-frame prediction method in accordance with an embodiment of the present invention discloses a method of determining the reference picture index of a temporal merging candidate. An intra-prediction unit can include a merging candidate deriving unit for deriving the spatial merging candidate and temporal merging candidate blocks of a current block and deriving a spatial merging symbol from the spatial merging candidate block and a temporal merging candidate from the temporal merging candidate block. A method of deriving the merging candidates will be additionally described later. [0054] The reconstructed residual block and the prediction block are added by the adder 255 , and the added block can experience the filtering unit 260 . The filtering unit 260 can apply one or more of a deblocking filter, an SAO, and an ALF to a reconstructed block or a reconstructed picture. The filtering unit 260 can output the reconfigured picture. The reconstructed picture can be stored in the reference picture buffer 270 and used for inter-prediction. [0055] A method of improving the prediction performance of an image coder and an image decoder includes a method of increasing the accuracy of an interpolation image and a method of predicting a difference signal. The difference signal indicates a difference between the original image and a prediction image. In the present invention, a “difference signal” can be replaced with a “residual signal”, a “residual block”, or a “difference block” depending on context. A person having ordinary skill in the art can distinguish the residual signal, the residual block, and the difference block from each other within a range that does not affect the spirit and essence of the invention. [0056] In an embodiment of the present invention, a term, such as a Coding Unit (CU), a Prediction Unit (PU), or a Transform Unit (TU), can be used as a unit for processing an image. [0057] The CU is an image processing unit on which coding/decoding are performed. The CU can include information used to code or decode a coding block, that is, a block unit set of luminance samples or chrominance samples on which coding/decoding are performed, and the samples of the coding block. [0058] The PU is an image processing unit on which prediction is performed. The PU can include information used to predict a prediction block, that is, a block unit set of luminance samples or chrominance samples on which prediction is performed, and the samples of the prediction block. Here, a coding block can be classified into a plurality of prediction blocks. [0059] The TU is an image processing unit on which transform is performed. The TU can include information used to transform a transform block, that is, a block unit set of luminance samples or chrominance samples on which transform is performed, and the samples of the transform block. Here, a coding block can be classified into a plurality of transform blocks. [0060] In an embodiment of the present invention, a block and a unit can be interpreted as having the same meaning unless described otherwise hereinafter. [0061] Furthermore, a current block can designate a block on which specific image processing is being performed, such as a prediction block on which prediction is now performed or a coding block on which prediction is now performed. For example, if one coding block is partitioned into two prediction blocks, a block on which prediction is now performed, from among the partitioned prediction blocks, can be designated as a current block. [0062] In an embodiment of the present invention, an image coding method and an image decoding method to be described later can be performed by the elements of the image coder and image decoder described with reference to FIGS. 1 and 2 . The element can include not only a hardware meaning, but also a software processing unit that can be performed by an algorithm. [0063] Hereinafter, a method of setting the reference picture index of a temporal merging candidate disclosed in an embodiment of the present invention can be used both in SKIP mode in an image processing method and merge mode, that is, one of modes, in an inter-frame prediction method. SKIP mode is an image processing method of outputting a block, predicted based on motion prediction information derived from surrounding blocks, as a reconstructed block without generating a residual block. Merge mode, that is, one of modes, in an inter-frame prediction method is an image processing method which is the same as SKIP mode in that a block is predicted based on motion prediction information derived from surrounding blocks, but is different from SKIP mode in that a block reconstructed by adding a residual block and a prediction block by coding and decoding information on the residual block is outputted. Intra-loop filtering methods, such as deblocking filtering and a sample adaptive offset, can be additionally applied to the outputted reconstructed block. [0064] FIG. 3 is a conceptual diagram illustrating an inter-frame prediction method using merge mode in accordance with an embodiment of the present invention. [0065] Referring to FIG. 3 , the inter-frame prediction using merge mode can be performed as follows. [0066] The inter-frame prediction method using merge mode refers to a method of deriving a merging candidate from a block neighboring a current block and performing inter-frame prediction by using the derived merging candidate. The neighboring block used to derive the merging candidate can be partitioned into a block which is located in the same picture as a current block and neighbors the current block and a block which is located in a picture different from a picture including a current block and at a location collocated with the current block. [0067] Hereinafter, in an embodiment of the present invention, from among neighboring blocks used to derive a merging candidate, a block which is located in the same picture as a current block and neighbors the current block is defined as a spatial merging candidate block, and motion prediction-related information derived from the spatial merging candidate block is defined as a spatial merging candidate. Furthermore, from among neighboring blocks used to derive a merging candidate, a block which is located in a picture different from a picture including a current block and at a location collocated with the current block is defined as a temporal merging candidate block, and motion prediction-related information derived from the temporal merging candidate block is defined as a temporal merging candidate. [0068] That is, the inter-frame prediction method using merge mode is an inter-frame prediction method for predicting a current block by using motion prediction-related information (i.e., a spatial merging candidate) on a spatial merging candidate block or motion prediction-related information (i.e., a temporal merging candidate) on a temporal merging candidate block to be described later. [0069] For example, motion vectors mvL0/L1, reference picture indices refldxL0/L1, and pieces of reference picture list utilization information predFlagL0/L1 can be used as the motion prediction-related information. [0070] FIG. 3(A) shows the motion vectors mvL0/L1, the reference picture indices refldxL0/L1, and the pieces of reference picture list utilization information predFlagL0/L1. [0071] A motion vector 304 is directional information and can be used for a prediction block to derive information on a pixel, located at a specific location, from a reference picture in performing inter-frame prediction. If inter-frame prediction is performed using a plurality of pieces of directional information in a prediction block, motion vectors for respective directions can be indicated by mvL0/L1. [0072] A reference picture index 306 is information on the index of a picture to which a prediction block refers in performing inter-frame prediction. If inter-frame prediction is performed using a plurality of reference pictures, reference pictures can be indexed using respective reference picture indices refldxL0 and refldxL1. [0073] The reference picture list utilization information can indicate that a reference picture has been derived from what reference picture list 0 308 . For example, pictures i, j, and k can be stored in a reference picture list 0 308 and used. If there are two lists in which a reference picture is stored, information on that the reference picture has been derived from what reference picture list can be indicated by predFlagL0 and predFlagL1. [0074] In order to perform the inter-frame prediction method using merge mode, first, a spatial merging candidate can be obtained through the following step (1). FIG. 3(B) discloses a spatial merging candidate and a temporal merging candidate. [0075] (1) A spatial merging candidate is derived from neighboring blocks for a current block (i.e., a target prediction block). [0076] As described above, a spatial merging candidate is motion prediction-related information derived from a spatial merging candidate block. The spatial merging candidate block can be derived on the basis of the location of a current block. [0077] Referring to FIG. 3(B) , the existing spatial merging candidate blocks 300 , 310 , 320 , 330 , and 340 have been derived based on a target prediction block. It is assumed that the location of a pixel present at an upper left end of the target prediction block is (xP, yP), the width of a prediction block is nPbW, the height of the target prediction block is nPbH, and MinPbSize is the smallest size of the prediction block. In an embodiment of the present invention hereinafter, the spatial merging candidate blocks of the prediction block can include a block including a pixel present at (xP−1, yP+nPbH), that is, a first block (or an A0 block) 300 on the left side, a block including a pixel present at (xP−1, yP+nPbH−1), that is, a second block (or an A1 block) 310 on the left side, a block including a pixel present at (xP+nPbW, yP−1), that is, a first block (or a B0 block) 320 at the upper end, a block including a pixel present at (xP+nPbW−1, yP−1), that is, a second block (or a B1 block) 330 at the upper end, and a block including a pixel present at (xP−1, yP−1), that is, a third block (or a B2 block) 340 at the upper end. Another value, for example, “MinPbSize” may be used instead of 1. In this case, a block at the same location can be indicated. Coordinates used to indicate the block at the specific location are arbitrary, and the block at the same location may be indicated by various other representation methods. [0078] The locations of the spatial merging candidate blocks 300 , 310 , 320 , 330 , and 340 and the number thereof and the locations of the temporal merging candidate blocks 360 and 370 and the number thereof disclosed in FIG. 3 are illustrative, and the locations of spatial merging candidate blocks and the number thereof and the locations of temporal merging candidate blocks and the number thereof can be changed if they fall within the essence of the present invention. Furthermore, order of merging candidate blocks preferentially scanned when a merging candidate list is configured may be changed. That is, the locations of candidate prediction blocks, the number thereof, and a scan order thereof, and a candidate prediction group used when a candidate prediction motion vector list is configured, described in the following embodiment of the present invention, are only illustrative and can be change if they fall within the essence of the present invention. [0079] A spatial merging candidate can be derived from an available spatial merging candidate block by determining whether the spatial merging candidate blocks 300 , 310 , 320 , 330 , and 340 are available or not. Information indicating whether a spatial merging candidate can be derived from a spatial merging candidate block or not is availability information. For example, if a spatial merging candidate block is located outside a slice, tile, or a picture to which a current block belongs or is a block on which intra-frame prediction has been performed, a spatial merging candidate, that is, motion prediction-related information, cannot be derived from the corresponding block. In this case, the spatial merging candidate block can be determined to be not available. In order to determine availability information on the spatial merging candidate, some determination methods can be used and embodiments thereof are described in detail later. [0080] If a spatial merging candidate block is available, motion prediction-related information can be derived and used to perform inter-frame prediction using merge mode on a current block. [0081] One coding block can be partitioned into one or more prediction blocks. That is, a coding block can include one or more prediction blocks. If a plurality of prediction blocks is included in a coding block, each of the prediction blocks can be indicated by specific index information. For example, if one coding block is partitioned into two prediction blocks, the two prediction blocks can be indicated by setting the partition index of one prediction block to 0 and the partition index of the other prediction block to 1. If a partition index is 0, a prediction block may be defined as another term, such as a first prediction block. If a partition index is 1, a prediction block may be defined as another term, such as a second prediction block. If one coding block is further partitioned into additional prediction blocks, index values indicative of the prediction blocks can be increased. The terms defined to designate the prediction blocks are arbitrary, and the terms may be differently used or differently interpreted. The partition index of a prediction block may also be used as information indicative of order that image processing, such as coding and decoding, is performed when a prediction block performs the image processing. [0082] If a plurality of prediction blocks is present within one coding block, there may be a case where coding or decoding on another prediction block must be first performed when a spatial merging candidate for the prediction block is derived. In accordance with an embodiment of the present invention, a method of deriving spatial merging candidates and temporal merging candidates in parallel to each of prediction blocks included in one coding block when generating a merging candidate list is additionally disclosed in detail. [0083] (2) Determine the reference picture index of a temporal merging candidate. [0084] A temporal merging candidate is motion prediction-related information derived from a temporal merging candidate block that is present at a picture different from a picture including a current block. The temporal merging candidate block is derived based on a block that is at a location collocated based on the location of the current block. The term ‘colocated block’ can be used as the same meaning as the temporal merging candidate block. [0085] Referring back to FIG. 3 , the temporal merging candidate blocks 360 and 370 can include the block 360 including a pixel at a location (xP+nPSW, yP+nPSH) in the colocated picture of a current prediction block or the block 370 including a pixel at a location (xP+(nPSW>>1), yP+(nPSH>>1)) if the block 360 including the pixel at the location (xP+nPSW, yP+nPSH) is not available, on the basis of the pixel location (xP, yP) within the picture including the prediction block. The prediction block 360 including the pixel at the location (xP+nPSW, yP+nPSH) in the colocated picture can be called a first temporal merging candidate block (or a first colocated block) 360 , and the prediction block including the pixel at the location (xP+(nPSW>>1), yP+(nPSH>>1)) in the colocated picture can be called a second temporal merging candidate block 370 . [0086] Finally, the final temporal merging candidate block used to derive a temporal merging candidate (or motion prediction-related information) can be at a location partially moved on the basis of the locations of the first temporal merging candidate block 360 and the second temporal merging candidate block 370 . For example, if only pieces of motion prediction-related information on some prediction blocks present in a colocated picture are stored in memory, a block at a location partially moved on the basis of the locations of the first temporal merging candidate block 360 and the second temporal merging candidate block 370 can be used as the final temporal merging candidate block for deriving the final motion prediction-related information. Like in a spatial merging candidate block, the location of a temporal merging candidate block can be changed or added unlike in FIG. 3 , and an embodiment thereof is described later. [0087] The reference picture index of a temporal merging candidate is information indicative of a picture that is referred in order for a current block to perform inter-frame predict on the basis of a motion vector mvLXCol derived from a temporal merging candidate. [0088] FIG. 4 is a conceptual diagram illustrating inter-frame prediction using a temporal merging candidate and the reference picture index of the temporal merging candidate in accordance with an embodiment of the present invention. [0089] Referring to FIG. 4 , a current block 400 , a picture 410 including the current block, a temporal merging candidate block (or a colocated block) 420 , and a colocated picture 430 including the colocated block can be defined. [0090] From a viewpoint of the temporal merging candidate block 420 , there is a picture 440 used in inter-frame prediction by the temporal merging candidate block in order to perform the inter-frame prediction on the temporal merging candidate block 420 . This picture is defined as the reference picture 440 of the colocated picture 430 . Furthermore, a motion vector that is used by the temporal merging candidate block 420 in order to perform inter-frame prediction from the reference picture 440 of the colocated picture 430 can be defined as mvCol 470 . [0091] From a standpoint of the current block 400 , a reference picture 460 used in the inter-frame prediction of the current block 400 on the basis of the calculated mvCol 470 has to be defined. The reference picture defined to be used in the inter-frame prediction of the current block 400 can be called the reference picture 460 of a temporal merging candidate. That is, the index of the reference picture 460 of the temporal merging candidate (i.e., the reference index of the temporal merging candidate) is a value indicative of a reference picture used in the temporal motion prediction of the current block 400 . At the step (2), the reference picture index of a temporal merging candidate can be determined. [0092] A mvCol 470 , that is, a motion vector derived from the temporal merging candidate block 420 , can be scaled and transformed into a different value depending on the distance between the colocated picture 430 and the reference picture 440 of the colocated picture and the distance between the picture 410 including the current block and the reference picture 460 of the temporal merging candidate derived through the step (2). [0093] That is, inter-frame prediction according to the temporal merging candidate of the current block 400 can be performed based on mvLXCol 480 derived through a step (3) to be described later, on the basis of the reference picture index 460 of the temporal merging candidate derived through the step (2) and the reference picture index 460 of the temporal merging candidate. mvLXCol can be defined as a temporal motion vector. [0094] In the existing image coding/decoding methods, the reference picture index of a temporal merging candidate can be determined based on the reference picture index candidate of a temporal merging candidate derived from the reference picture index of a spatial merging candidate in a target prediction block. If this method is used, there may be a case where the reference picture index of a spatial merging candidate that has not yet been coded or decoded must be derived. In this case, the reference picture index of the spatial merging candidate can be derived only when coding or decoding on a prediction block including the corresponding spatial merging candidate is finished. Accordingly, if the reference picture index of a temporal merging candidate is determined based on the reference picture index candidates of temporal merging candidates derived from all spatial merging candidate blocks, a process of deriving the reference pictures of the temporal merging candidates for a current block cannot be performed in parallel. FIG. 5 discloses this problem. [0095] FIG. 5 is a conceptual diagram illustrating a case where one coding block is partitioned into two prediction blocks. [0096] Referring to FIG. 5 , one coding block is partitioned into a first prediction block 500 and a second prediction block 520 having an N×2N form. Spatial merging candidate blocks for the first prediction block 500 are derived on the basis of the location of the first prediction block 500 as in FIG. 5(A) , and spatial merging candidate blocks for the second prediction block 520 are derived on the basis of the location of the second prediction block 520 as in FIG. 5(B) . Although not shown, in temporal merging candidate blocks, temporal merging candidates can be derived on the basis of the location of each of prediction blocks. [0097] The spatial merging candidate blocks of the first prediction block 500 are outside the first prediction block 500 and are at locations included in blocks on which coding or decoding has already been performed. [0098] In contrast, an A1 block 530 , from among the spatial merging candidate blocks of the second prediction block 520 , is present within the first prediction block 500 . Accordingly, after prediction on the first prediction block 500 is performed, motion prediction-related information (e.g., a motion vector, a reference picture index, and reference picture list utilization information) on the A1 block 530 can be known. Furthermore, the motion prediction-related information of the A0 block 550 cannot be derived because the A0 block 550 is at a location that has not yet been coded or decoded. [0099] If the reference picture index of a temporal merging candidate is derived from the motion prediction-related information of the A1 block 530 , it can be derived after coding and decoding on the first prediction block 500 are finished. Furthermore, the reference picture index cannot be derived from the A0 block 550 . That is, since the reference picture indices of some spatial merging candidate blocks cannot be derived, the reference picture indices of temporal merging candidates for respective prediction blocks cannot be derived in parallel. [0100] In an embodiment of the present invention, in order to solve the problem, methods of deriving the reference picture indices of temporal merging candidates (or the reference indices of temporal merging candidates) for prediction blocks are disclosed. [0101] If a method of deriving the reference picture indices of temporal merging candidates in accordance with an embodiment of the present invention is used, processes of deriving the reference picture indices of temporal merging candidates for some prediction blocks can be performed in parallel. Since the reference picture indices of temporal merging candidates are derived in parallel, inter-frame prediction processes using merge mode for a plurality of prediction blocks included in one coding block can be performed in parallel. [0102] Hereinafter, in an embodiment of the present invention, a method of deriving the reference picture index of a temporal merging candidate is disclosed and additionally described in detail. [0103] (3) Derive motion prediction-related information on a temporal merging candidate block. [0104] At the step (3), in order to perform motion prediction based on a temporal merging candidate, temporal merging candidates, such as information on whether a temporal merging candidate block is available or not (availableFlagCol), reference picture list utilization information (PredFlagLXCol), and information on the motion vector (mvLXCol) of a temporal merging candidate, can be derived. The motion prediction-related information derived from the temporal merging candidate can be defined as a term ‘temporal merging candidate’. The availability information on the temporal merging candidate block indicates whether a temporal merging candidate can be derived from the temporal merging candidate block or not. The temporal merging candidate can be included in a merging candidate list on the basis of the availability information on the temporal merging candidate block. [0105] (4) Derive a merging candidate list. [0106] A merging candidate list can include information on a merging candidate that can be used in inter-frame prediction using merge mode on the basis of availability information on a merging candidate block (i.e., a spatial merging candidate block or a temporal merging candidate block). One of merging candidates included in a merging candidate list can be used to perform inter-frame prediction using merge mode on a current block. Information on whether what merging candidate will be used to predict a current block (i.e., a merging index) can be coded in a coding step and transmitted to a decoder. [0107] A merging candidate list can be generated with the following order of priority. [0108] 1) If an A1 block is available, a merging candidate derived from the A1 block [0109] 2) If a B1 block is available, a merging candidate derived from the B1 block [0110] 3) If a B0 block is available, a merging candidate derived from the B0 block [0111] 4) If an A0 block is available, a merging candidate derived from the A0 block [0112] 5) If a B2 block is available, a merging candidate derived from the B2 block [0113] 6) If a Col block is available, a merging candidate derived from the Col block [0114] The merging candidate list can include, for example, 0 to 5 merging candidates depending on the number of available blocks. If the number of blocks used to derive a merging candidate is many, more merging candidates may be included in the merging candidate list. [0115] (5) If the number of merging candidates included in a merging candidate list is smaller than a maximum number of merging candidates that can be included in the merging candidate list, an additional merging candidate is derived. [0116] An additional merging candidate can be a candidate generated by combining pieces of motion prediction-related information on the existing merging candidates (i.e., a combined merging candidate) or can be a 0-vector merging candidate (i.e., a zero merging candidate). Here, the 0-vector merging candidate designates a merging candidate having a motion vector (0,0). [0117] (6) Determine a merging candidate applied to inter-frame prediction performed on a current block, from among merging candidates included in a merging candidate list, and set motion prediction-related information on the determined merging candidate as motion prediction-related information on a current block. [0118] In a decoding process, inter-frame prediction using merge mode can be performed on a current block on the basis of a merging index merge_idx[xP][yP], that is, information on which one of candidates included in a merging candidate list is used in inter-frame prediction performed on the current block. [0119] Through a procedure of the step (1) to the step (6), motion prediction-related information on a current block can be derived and inter-frame prediction can be performed on the current block based on the derived motion prediction-related information. [0120] An embodiment of the present invention discloses a method of deriving the reference picture indices of temporal merging candidates for a plurality of prediction blocks, included in one coding block, in parallel in setting the reference picture index of a temporal merging candidate at the step (2) is disclosed. [0121] Various kinds of methods below can be used as the method of deriving the reference picture indices of temporal merging candidates for a plurality of prediction blocks, included in a coding block, in parallel. [0122] 1) A method of setting the location of a spatial merging candidate block, used to derive the reference picture index candidate of a temporal merging candidate for a target prediction block (i.e., a current block), as a location at which a coding block including the current block is located and on which coding or decoding has already been performed. [0123] 2) A method of, if the location of a spatial merging candidate block used to derive the reference picture index candidate of a temporal merging candidate for a target prediction block (i.e., a current block) is within a coding block or a location on which coding has not yet been performed, setting the reference picture index candidate of a temporal merging candidate derived from a spatial merging candidate at the corresponding location to ‘0’. [0124] 3) A method of setting the reference picture index of the temporal merging candidate of a target prediction block (i.e., a current block) to ‘0’ that is a fixed value. [0125] 4) A method of, if the location of a spatial merging candidate block referred to derive the reference picture index candidate of the temporal merging candidate of a target prediction block (i.e., e current block) is within a coding block or a location on which coding has not yet been performed, not using the reference picture index of the spatial merging candidate block at the corresponding location in order to derive the reference picture index of the temporal merging candidate. [0126] 5) A method of previously determining a spatial merging candidate block at a specific location that is referred to derive the reference picture index of the temporal merging candidate of a target prediction block (i.e., a current block) and deriving the reference picture index of the temporal merging candidate from the spatial merging candidate block at the specific location. [0127] 6) A method of, if the locations of some of the spatial merging candidate blocks of spatial merging candidates derived to perform mergence on a target prediction block (i.e., a current block) are within a coding block or locations on which coding has not yet been performed and thus pieces of information on the reference picture indices of temporal merging candidates cannot be derived from the spatial merging candidate blocks at the corresponding locations, fixing the spatial merging candidate blocks at the corresponding locations as locations outside the coding block on which coding or decoding has been performed. [0128] The following embodiments of the present invention disclose the methods of deriving the reference picture index of a temporal merging candidate in detail. [0129] First, problems occurring when determining the reference picture index of a temporal merging candidate in the prior art, described with reference to FIG. 5 , are described in detail with reference to FIG. 6 . [0130] FIG. 6 is a conceptual diagram illustrating a method of setting the reference picture index of a temporal merging candidate in accordance with an embodiment of the present invention. [0131] Referring to FIG. 6 , one coding block (e.g., a 2N×2N form) can be partitioned into two prediction blocks (e.g., N×2N). In the first prediction block 600 of the two partitioned blocks, all spatial merging candidate blocks 605 , 610 , 615 , 620 , and 625 are present outside the coding block. In contrast, in the second prediction block 650 of the two partitioned blocks, some (e.g., 655 , 665 , 670 , and 675 ) of spatial merging candidate blocks 655 , 660 , 665 , 670 , and 675 are present outside the coding block, and some (e.g., 660 ) of the spatial merging candidate blocks 655 , 660 , 665 , 670 , and 675 are present within the coding block. [0132] The reference picture index of a temporal merging candidate for a current block (i.e., a target prediction block) can be derived from the reference picture index of a spatial merging candidate. That is, the reference picture index of a temporal merging candidate for a current block can be derived based on information on a reference picture index that has been used by a spatial merging candidate block to perform inter-frame prediction. [0133] For example, it can be assumed that the reference picture indices of three of a plurality of spatial merging candidates for a current block are refldxLXA, refldxLXB, and refldxLXC. Pieces of information on the reference picture indices refldxLXA, refldxLXB, and refldxLXC can become the reference picture index candidates of temporal merging candidates, and the reference picture index values of the temporal merging candidates can be derived based on the reference picture index candidates of the temporal merging candidates. [0134] If the above method is used, spatial merging candidate blocks for a current block need to be coded or decoded in advance because pieces of information on the reference picture indices of the spatial merging candidate blocks for the current block are necessary to derive the reference picture indices of temporal merging candidates for the current block. [0135] Referring back to FIG. 6 , the first prediction block 600 is a block in which the spatial merging candidates are included in locations outside the coding block on which coding or decoding has already been performed as described above. Accordingly, if the first prediction block 600 is a current block on which prediction is performed, the reference picture index candidates of temporal merging candidates for the first prediction block 600 can be directly derived from the spatial merging candidate blocks of the first prediction block 600 . [0136] In the second prediction block 650 , however, some (e.g., 660 ) of the spatial merging candidates are present in the first prediction block 600 that is within the coding block as described above. Accordingly, when inter-frame prediction using merge mode is performed on the second prediction block 650 , the reference picture indices of temporal merging candidates for the first prediction block 650 cannot be derived until the A1 block 660 is coded or decoded, that is, until prediction is performed on the first prediction block 600 including the A1 block 660 . In this case, there is a problem in that inter-frame prediction using merge mode cannot be performed on the first prediction block 600 and the second prediction block 650 in parallel because the temporal merging candidates of the second prediction block 650 are not derived until prediction is performed on the first prediction block 600 . In order to solve the problem, a variety of methods can be used. [0137] Only some of the partition forms of a prediction block are disclosed in the following embodiments of the present invention, for convenience of description, but the present invention can be applied to the partition forms of several prediction blocks of a coding block and embodiments thereof are also included in the scope of the present invention. [0138] FIG. 7 is a conceptual diagram illustrating a method of deriving the reference picture indices of temporal merging candidates in accordance with an embodiment of the present invention. [0139] The embodiment of FIG. 7 discloses a method of setting the locations of spatial merging candidate blocks to which reference is made by a prediction block in order to derive the reference picture indices of temporal merging candidates as locations outside a coding block including a current prediction block. [0140] FIG. 7(A) shows a case where one coding block is partitioned into two prediction blocks 700 and 750 having an N×2N form. [0141] All the spatial merging candidate blocks of the first prediction block 700 are at locations outside a coding unit on which coding or decoding has already been performed. Thus, the reference picture index candidates of temporal merging candidates for the first prediction block 700 can be directly derived by using the already coded or decoded spatial merging candidate blocks. [0142] In the case of the second prediction block 750 , however, the locations of some (e.g., 710 and 720 ) of spatial merging candidate blocks used to derive the reference picture indices of temporal merging candidates can be changed, and the reference picture indices of the temporal merging candidates can be derived from the changed locations. [0143] In order to derive the reference picture index candidates of the temporal merging candidates, the spatial merging candidate block 710 can be replaced with a block 715 outside the coding block without using the spatial merging candidate block 710 included in the coding unit, from among the spatial merging candidate blocks of the second prediction block 750 , and the reference picture index of the block 715 can be used as the reference picture index candidate of a temporal merging candidate. [0144] Furthermore, the spatial merging candidate block 720 can be replaced with a block 725 outside the coding block without using the block 720 outside the coding unit on which coding or decoding has not yet been performed, from among the spatial merging candidate blocks, and the reference picture index of the block 725 can be used as the reference picture index candidate of a temporal merging candidate. [0145] That is, the reference picture index candidates of the temporal merging candidates can be derived by using an A0′ block 725 and an A1′ block 715 outside the coding block instead of the A0 block 710 and the A1 block 720 of the second prediction block 750 . [0146] If the above method is used, all the spatial merging candidate blocks used to derive the reference picture indices of the temporal merging candidates can become blocks included in an already coded block in the second prediction block 750 . Accordingly, in the second prediction block 750 , the reference picture indices of the temporal merging candidates can be derived irrespective of whether or not a prediction process has been performed on the first prediction block 700 . [0147] FIG. 7(B) shows a case where one coding block is partitioned into two prediction blocks having a 2N×N form. [0148] As in FIG. 7(A) , in FIG. 7(B) , instead of a B1 block 780 , that is, a block included within the coding block, and a B0 block 790 , that is, a block on which coding or decoding has not yet been performed, from among the spatial merging candidate blocks of a second prediction block 770 , a B1′ block 785 and a B0′ block 795 that are already coded blocks can be used to derive the reference picture indices of temporal merging candidates for the second prediction block 770 . [0149] FIG. 8 is a conceptual diagram illustrating a method of deriving the reference picture indices of temporal merging candidates in accordance with an embodiment of the present invention. [0150] The embodiment of FIG. 8 discloses a method of setting the reference picture index candidates of temporal merging candidates, derived from a spatial merging candidate block present within a coding block and a spatial merging candidate block present in a location on which coding or decoding has not yet been performed, to ‘0’, if the locations of spatial merging candidate blocks referred to derive the reference picture indices of temporal merging candidates for a target prediction block (i.e., a current block) are within the coding block including the current block or are locations on which coding or decoding has not yet been performed. [0151] FIG. 8(A) shows a case where one coding block is partitioned into two prediction blocks having an N×2N form. [0152] Referring to FIG. 8(A) , all the spatial merging candidate blocks of a first prediction block 800 are at locations outside the coding unit on which coding or decoding has already been performed. Accordingly, the reference picture index candidates of temporal merging candidates for the first prediction block 800 can be derived from the spatial merging candidate blocks of the first prediction block 800 . [0153] In the case of a second prediction block 850 , assuming that the reference picture indices of some spatial merging candidate blocks (e.g., 810 and 820 ) are ‘0’, the reference picture index candidates of temporal merging candidate for the second prediction block 850 can be derived. In relation to a spatial merging candidate block located within the coding block including a target prediction block (i.e., a current block) or a spatial merging candidate block at a location on which coding or decoding has not yet been performed, the reference picture index candidate of a temporal merging candidate derived from a corresponding spatial merging candidate block can be set to ‘0’ and the reference picture index of a temporal merging candidate for the current block can be derived from the set reference picture index candidate. [0154] For example, a process of setting the reference picture index candidates of temporal merging candidates, derived from the A0 block 810 and the A1 block 820 of the second prediction block 850 , to ‘0’ in advance when deriving the reference picture index candidates of the temporal merging candidates and deriving the reference picture indices of the temporal merging candidates from the set reference picture index candidates can be used. [0155] FIG. 8(B) shows a case where one coding block is partitioned into two prediction blocks having a 2N×N form. [0156] All the spatial merging candidate blocks of a first prediction block 860 are at locations outside the coding unit on which coding or decoding has been completed. Accordingly, the reference picture index candidates of temporal merging candidates for the first prediction block 860 can be directly derived from the spatial merging candidate blocks of the first prediction block 860 . [0157] The reference picture index candidate of a temporal merging candidate derived from a spatial merging candidate block 880 included in a prediction block on which prediction has not yet been performed or some spatial merging candidate blocks (e.g., 890 ) at locations on which a coding or decoding process has not yet been performed can be set to ‘0’, when deriving the reference picture indices of temporal merging candidates for a second prediction block 870 . The reference picture index candidates of the temporal merging candidates can be derived from the set reference picture index candidates. [0158] For example, the above method can be used in a process of setting the reference picture indices of temporal merging candidates derived from a B0 block 880 and a B1 block 890 , that is, the spatial merging candidate blocks of the second prediction block 870 , to ‘0’ and deriving the reference picture indices of the temporal merging candidates for the second prediction block 870 . [0159] FIG. 9 is a conceptual diagram illustrating a method of deriving the reference picture indices of temporal merging candidates in accordance with an embodiment of the present invention. [0160] The embodiment of FIG. 9 discloses a method in which a prediction block sets the reference picture index of a temporal merging candidate to ‘0’, that is, a fixed value, and using the set reference picture index. [0161] FIG. 9(A) shows a case where one coding block is partitioned into two prediction blocks having an N×2N form. [0162] Referring to FIG. 9(A) , in order to derive the reference picture indices of temporal merging candidates for a first prediction block 900 and a second prediction block 950 , the reference picture index values of temporal merging candidates can be set to ‘0’ and used, without using spatial merging candidate blocks 905 to 925 and 930 to 947 . If this method is used, the degree of complexity in the deriving of coding and decoding can be reduced and the speed of coding and decoding can be increased because a step of deriving the reference picture indices of temporal merging candidates is not performed. Furthermore, the reference picture indices of temporal merging candidates for a current block can be derived without a need to wait for until prediction on other prediction blocks included in a current coding block is performed. Accordingly, the reference picture indices of temporal merging candidates for a plurality of prediction blocks included in one coding block can be derived in parallel. [0163] FIG. 9(B) shows a case where one coding block is partitioned into two prediction blocks having a 2N×N form. [0164] Likewise, in FIG. 9(B) , in order to derive the reference picture indices of temporal merging candidates for a first prediction block 960 and a second prediction block 990 , the reference picture index values of the temporal merging candidates can be fixed to ‘0’ and used without using spatial merging candidates. [0165] In FIG. 9 , ‘0’ is marked in the spatial merging candidate blocks, for convenience of description. However, when actually deriving the reference picture indices of temporal merging candidates, a value set to ‘0’ can be used without a procedure for searching for the reference picture indices of the spatial merging candidate blocks. ‘0’ is only an example of a fixed picture index, and another picture index other than 0 may be used and embodiments thereof are also included in the scope of the present invention. [0166] FIG. 10 is a conceptual diagram illustrating a method of deriving the reference picture indices of temporal merging candidates in accordance with an embodiment of the present invention. [0167] The embodiment of FIG. 10 discloses a method of, if the location of a spatial merging candidate block referred to derive the reference picture index of a temporal merging candidate for a current block (i.e., a target prediction block) is within a coding block including the current block or at a location on which coding has not yet been performed, not using the reference picture index of the spatial merging candidate block as a candidate for deriving the reference picture index of the temporal merging candidate. [0168] FIG. 10(A) shows a case where one coding block is partitioned into two prediction blocks having an N×2N form. [0169] Referring to FIG. 10(A) , the A1 block 1030 and the A0 block 1020 of a second prediction block 1010 are a block within the coding block including a current block and a block at a location on which coding or decoding has not yet been performed. Pieces of information on the reference picture indices of the A1 block 1030 and the A0 block 1020 cannot be used when deriving the reference picture indices of temporal merging candidates for a first prediction block 1000 . [0170] Accordingly, when deriving the reference picture indices of the temporal merging candidates from the second prediction block 1010 , the pieces of information on the reference picture indices of the A1 block 1030 and the A0 block 1020 can be set to ‘−1’. If the reference picture index value of a specific spatial merging candidate block is ‘−1’, the spatial merging candidate block can indicate a block that is not used to derive the reference picture index of a temporal merging candidate. [0171] FIG. 10(B) shows a case where one coding block is partitioned into two prediction blocks having a 2N×N form. [0172] Referring to FIG. 10(B) , the B1 block 1060 of a second prediction block 1050 is a spatial merging candidate block within the coding block and is a block whose reference picture index information can be known only when prediction is performed on a first prediction block 1040 . The B0 block 1070 of the second prediction block 1050 is a spatial merging candidate block at a location on which coding has not yet been performed, and information on the reference picture index thereof cannot be known. [0173] In this case, in order to derive the reference picture indices of temporal merging candidates from the first prediction block 1040 and the second prediction block 1050 in parallel, pieces of information on the reference picture indices of the B1 block 1060 and the B0 block 1070 can be set to ‘−1’. That is, the B0 block 1070 and the B1 block 1060 may not be used as blocks for deriving the reference picture index candidates of temporal merging candidates for the second prediction block 1050 . [0174] FIG. 11 is a conceptual diagram illustrating a method of deriving the reference picture index of a temporal merging candidate in accordance with an embodiment of the present invention. [0175] The embodiment of FIG. 11 discloses a method of previously determining a specific spatial merging candidate block referred by a prediction block in order to derive the reference picture index of a temporal merging candidate and deriving the reference picture index of the temporal merging candidate from the specific spatial merging candidate block. [0176] FIG. 11(A) shows a case where one coding block is partitioned into two prediction blocks having an N×2N form. [0177] A first prediction block 1100 and a second prediction block 1120 can share spatial merging candidate blocks A0, A1, B0, B1, and B2. That is, the spatial merging candidate blocks A0, A1, B0, B1, and B2 used to perform inter-frame prediction using merge mode in the first prediction block 1100 and the second prediction block 1120 can be blocks outside the coding block. [0178] A reference picture index for the temporal merging of the first prediction block 1100 can be set to the reference picture index value of the B1 block 1105 . That is, the fixed reference picture index of a spatial merging candidate block at a specific location of a prediction block can be set to a reference picture index value for the temporal merging of a current block depending on a partition form. [0179] If the B1 block 1125 is not available, the reference picture index value can be set to ‘0’ and used. [0180] Like in the second prediction block 1120 , the reference picture index value of the A1 block 1125 can be used as a reference picture index for temporal merging. If the B1 block 1105 is not available, the reference picture index value can be set to ‘0’ and used. [0181] FIG. 11(B) shows a case where one coding block is partitioned into two prediction blocks having a 2N×N form. [0182] A first prediction block 1150 and a second prediction block 1170 can share spatial merging candidate blocks A0, A1, B0, B1, and B2. That is, spatial merging candidate blocks for performing inter-frame prediction using merge mode in the first prediction block 1150 and the second prediction block 1170 can be blocks outside the coding block. [0183] A reference picture index for the temporal merging of the first prediction block 1150 can be set to the reference picture index value of the A1 block 1155 . That is, the reference picture index of a spatial merging candidate block at a specific location of a prediction block can be set to a reference picture index value for the temporal merging of a current block depending on a partition form. [0184] If the B1 block 1175 is not available, the reference picture index value can be set to ‘0’ and used. [0185] Like in the second prediction block 1170 , the reference picture index value of the B1 block 1175 can be used as a reference picture index for temporal merging. If the B1 block 1175 is not available, the reference picture index value can be set to ‘0’ and used. [0186] FIG. 12 is a conceptual diagram illustrating a method of deriving the reference picture index of a temporal merging candidate in accordance with an embodiment of the present invention. [0187] The embodiment of FIG. 12 discloses a method of previously determining a specific spatial merging candidate block referred by a target prediction block in order to derive the reference picture index of a temporal merging candidate and deriving the reference picture index of the temporal merging candidate from the specific spatial merging candidate block. [0188] Referring to FIG. 12 , different spatial merging candidate blocks can be used to derive the reference picture index of a temporal merging candidate depending on a form of a prediction block partitioned from one coding block. [0189] For example, in a prediction block, one of an A1 block and a B1 block, from among spatial merging candidate blocks, can be used as a block for deriving the reference picture index of a temporal merging candidate. From among the two spatial merging candidate blocks, a spatial merging candidate block within a coding block is not used to derive the reference picture index of the temporal merging candidate, but a spatial merging candidate block outside the coding block can be used to derive the reference picture index of the temporal merging candidate. [0190] Although FIG. 12(A) showing a case where a coding block is partitioned into prediction blocks having an N×2N form and FIG. 12(B) showing a case where a coding block is partitioned into prediction blocks having a 2N×N form are illustrated for convenience of description, the same method can be applied to a coding block partitioned in various forms. [0191] FIG. 12(A) shows a case where a coding block is partitioned into prediction blocks having an N×2N form. [0192] If one coding block is partitioned into prediction blocks having an N×2N form, the reference picture index of a B1 block 1220 , that is, a spatial merging candidate located outside the coding block and at a location on which coding or decoding has already been performed, from among two spatial merging candidate blocks (e.g., an A1 block 1200 and the B1 block 1220 ), can be set as the reference picture index of a temporal merging candidate for a second prediction block 1210 . [0193] FIG. 12(B) shows a case where a coding block is partitioned into prediction blocks having a 2N×N size. [0194] If one coding block is partitioned into prediction blocks having a 2N×N form, the reference picture index of an A1 block 1240 , that is, a spatial merging candidate outside the coding block, from among two spatial merging candidate blocks (e.g., the A1 block 1240 and a B1 block 1260 ), can be set as the reference picture index of a temporal merging candidate for a second prediction block 1250 . [0195] FIG. 13 is a conceptual diagram illustrating a method of deriving the reference picture index of a temporal merging candidate in accordance with an embodiment of the present invention. [0196] The embodiment of FIG. 13 discloses a method of, if the locations of some of the spatial merging candidates of a prediction block are within a coding block or placed at locations on which coding has not yet been performed, fixing the locations of the spatial merging candidate blocks of the corresponding prediction block to locations outside the coding block and using the fixed locations. [0197] FIG. 13(A) shows a case where one coding block is partitioned into two prediction blocks having an N×2N form. [0198] A first prediction block 1300 can determine spatial merging candidate blocks 1305 , 1310 , 1315 , 1320 , and 1325 on the basis of the first prediction block 1300 . In contrast, a second prediction block 1330 can fix spatial merging candidate blocks to blocks 1335 , 1340 , 1345 , 1350 , and 1355 placed at locations outside the coding block and use the fixed spatial merging candidate blocks. That is, the spatial merging candidate blocks 1335 , 1340 , 1345 , 1350 , and 1355 can be derived on the basis of the coding block, and the derived spatial merging candidate blocks 1335 , 1340 , 1345 , 1350 , and 1355 can be used in inter-frame prediction using merge mode for the second prediction block 1330 . [0199] FIG. 13(B) shows a case where one coding block is partitioned into two prediction blocks having a 2N×N form. [0200] Likewise, in FIG. 13(B) , a first prediction block can use spatial merging candidate blocks derived on the basis of a prediction block. In contrast, the spatial merging candidate blocks 1365 , 1370 , 1375 , 1380 , and 1385 of a second prediction block 1360 can be derived on the basis of the coding block. [0201] FIG. 14 is a flowchart illustrating a method of including a temporal merging candidate in a merging candidate list in accordance with an embodiment of the present invention. [0202] The embodiment of FIG. 14 discloses a process of deriving the reference picture index of a temporal merging candidate by using an index value calculated by the above-described method of deriving the reference picture index of a temporal merging candidate and of including the temporal merging candidate in a merging candidate list. [0203] Referring to FIG. 14 , the reference picture index of a temporal merging candidate is derived at step S 1400 . [0204] The reference picture index of the temporal merging candidate refers to the reference picture index of a picture referred by a current block (i.e., a target prediction block) in order to perform inter-frame prediction using merge mode as described above. The reference picture index of the temporal merging candidate can be derived by several methods of deriving the reference picture indices of temporal merging candidates in parallel in relation to a prediction block. For example, the reference picture index of the temporal merging candidate can be derived by several methods, such as 1) a method of always placing the spatial location of a spatial merging candidate block to be referred outside a coding block, 2) a method of replacing a reference picture index value, derived from a spatial merging candidate block to be referred, with ‘0’ if the spatial location of the spatial merging candidate block is within a coding block, and 3) a method of fixing the reference picture index of a temporal merging candidate to ‘0’ unconditionally. [0205] A temporal merging candidate is derived at step S 1410 . [0206] As described above, the temporal merging candidate can be motion prediction-related information (e.g., predFlag or mvLXCol) derived from a prediction block (e.g., a first temporal merging candidate block) which includes a pixel at a location (xP+nPbW, yP+nPbH) in the colocated picture of a current block on the basis of the location (xP, yP) of a pixel within a picture including the prediction block. If the prediction block including the pixel at the location (xP+nPbW, yP+nPbH) in the colocated picture is not available or is a block predicted by an intra-frame prediction method, motion prediction-related information (e.g., a temporal merging candidate) can be derived from a prediction block (e.g., a second temporal merging candidate block) including a pixel at a location (xP+(nPbW>>1), yP+(nPbH>>1)). [0207] Finally, a final temporal merging candidate block (i.e., a colocated block) used to derive the motion prediction-related information can be a block at a location that has been partially moved on the basis of the locations of the first temporal merging candidate block and the second temporal merging candidate block. For example, if only pieces of motion prediction-related information on some blocks are stored in memory, a temporal merging candidate block present at a location partially moved on the basis of the locations of the first temporal merging candidate block and the second temporal merging candidate block can be determined as the final colocated block for deriving the temporal merging candidate (i.e., motion prediction-related information). [0208] In deriving the temporal merging candidate, different temporal merging candidates can be derived depending on whether a current block is a block using a single merging candidate list or a block not using a single merging candidate list. If the current block is a block using a single merging candidate list, a plurality of prediction blocks included in a coding block can use temporal merging candidates derived from one temporal merging candidate block. If the current block is a block not using a single merging candidate list, a merging candidate list for a plurality of prediction blocks included in a coding block can be generated and inter-frame prediction using merge mode can be performed individually. That is, in this case, the inter-frame prediction can be performed by using temporal merging candidates derived from the temporal merging candidate block for each prediction block. An example in which inter-frame prediction is performed by using a single merging candidate list is described below. [0209] FIG. 15 is a conceptual diagram illustrating a method of generating a single merging candidate list by sharing all spatial merging candidates and temporal merging candidates in a plurality of prediction blocks in accordance with an embodiment of the present invention. [0210] The embodiment of FIG. 15 discloses a method of a plurality of prediction blocks, partitioned from one coding block, generating a single merging candidate list by sharing all spatial merging candidates and temporal merging candidates determined based on the coding block. [0211] Referring to FIG. 15 , a first prediction block 1500 and a second prediction block 1550 can derive spatial merging candidates from the same spatial merging candidate block and share the derived spatial merging candidates. Spatial merging candidate blocks for the first prediction block 1500 and the second prediction block 1550 are blocks determined based on a coding block, and an A0 block 1505 , an A1 block 1510 , a B0 block 1515 , a B1 block 1520 , and a B2 block 1525 can be used as the spatial merging candidate blocks. [0212] The location of each of the spatial merging candidate blocks can be a location including a pixel shown in the drawing on the basis of the upper left location (xC, yC) and nCS (i.e., the size of the coding block) of the coding block. [0213] The A0 block 1505 can be a block including a pixel at a location (xC−1, yC+nCS), the A1 block 1510 can be a block including a pixel at a location (xC−1, yC+nCS-1), the B0 block 1515 can be a block including a pixel at a location (xC+nCS, yC−1), the B1 block 1520 can be a block including a pixel at a location (xC+nCS-1, yC−1), and the B2 block 1525 can be a block including a pixel at a location (xC−1, yC−1). [0214] Furthermore, the first prediction block 1500 and the second prediction block 1550 can share temporal merging candidates. Temporal merging candidate blocks 1560 and 1570 for deriving the temporal merging candidates shared by the first prediction block 1500 and the second prediction block 1550 can be blocks at locations derived on the basis of the upper left locations (xC, yC) of the coding block and the size nCS of the coding block. [0215] The reference picture indices of the temporal merging candidates can be derived by the aforementioned methods. [0216] For example, the temporal merging candidate blocks 1560 and 1570 can include the prediction block 1560 including a pixel at a location (xC+nCS, yC+nCS) in the colocated picture of a current prediction block on the basis of the pixel location (xC, yC) within the picture including the prediction block or can be the prediction block 1570 including a pixel at a location (xC+(nCS>>1), yC+(nCS>>1)) if the prediction block including the pixel at the location (xC+nCS, yC+nCS) is not available. [0217] If temporal merging candidates are not shared, each of the temporal merging candidates for the first prediction block 1500 and the second prediction block 1550 can be derived. [0218] If a method of deriving a single merging candidate list is used, inter-frame prediction can be performed by parallel merging processing performed on each prediction block, and a merging candidate list for each prediction block does not need to be derived separately. Accordingly, by using a single merging candidate list in accordance with an embodiment of the present invention, an image processing speed can be improved in apparatuses, such as Ultra-High Definition TeleVision (UHDTV) that requires a large amount of data processing. [0219] FIG. 15 discloses only the first N×2N prediction block 1500 and the second N×2N prediction block 1550 each partitioned in an N×2N form, but this method can also be applied to prediction blocks partitioned in various forms, such as blocks having different partition forms (e.g., 2N×N, 2N×nU, 2N×nD, nL×2N, nR×2N, and N×N). [0220] Furthermore, in this method, whether or not to apply a single merging candidate list can be differently determined depending on the size of a block or a partition depth. For example, information on whether a single merging candidate list can be used in a specific block or not can be derived on the basis of pieces of information on the size of a block and the size of a coding block on which a merging process can be performed in parallel. For example, information on whether a single merging candidate list can be used in a specific block or not can be represented by flag information. A flag indicating whether or not a single merging candidate list can be used in a specific block can be defined as singleMCLflag (i.e., a single merge candidate list flag). For example, if the single merge candidate list flag singleMCLflag is 0, it can indicate that a block does not use a single merging candidate list. If the single merge candidate list flag singleMCLflag is 1, it can indicate that a block uses a single merging candidate list. Spatial merging candidates for a prediction block can be derived on the basis of a coding block based on a value of the single merge candidate list flag singleMCLflag. [0221] For example, the size of a block on which a merging process can be performed in parallel can derive flag information indicating that a prediction block, partitioned from an 8×8 coding block on the basis of information indicative of a value greater than a 4×4 size and information indicating that the size of a current block is 8×8, uses a single merging candidate list. The derived flag can be used to derive the spatial merging candidates and temporal merging candidates of a prediction block on the basis of a coding block. [0222] Referring back to FIG. 14 , availability information on the temporal merging candidate and a temporal motion vector can be derived on the basis of information on the reference picture index of the temporal merging candidate derived to derive the temporal merging candidate at the step S 1410 . [0223] The availability information on the temporal merging candidate can be used as information indicating whether the temporal merging candidate can be derived on the basis of a temporal merging candidate block. The temporal motion vector can be derived if the temporal merging candidate is available. [0224] Referring back to FIG. 4 , the temporal motion vector mvLXCol can be scaled and derived on the basis of the distance between two pictures derived based on the index of the picture 430 including a temporal merging candidate and the index of the reference picture 440 referred by the colocated picture 410 and the distance between pictures derived based on the index of the colocated picture 410 including the current block 400 and the index of the reference picture of a temporal merging candidate (i.e., the index of the reference picture 460 referred by the current block 400 in inter-frame prediction). [0225] If the temporal merging candidate is available, the temporal merging candidate is included in a merging candidate list at step S 1420 . [0226] When configuring the merging candidate list, if the temporal merging candidate is available based on availability information on the temporal merging candidate derived at the step S 1410 , a corresponding block can be included in the merging candidate list. [0227] FIG. 16 discloses a method in which prediction blocks within the same coding block share spatial merging candidates and temporal merging candidates only when the size of a block is equal to or smaller than a specific size. [0228] FIG. 16 is a conceptual diagram illustrating a method of generating a single candidate list in accordance with an embodiment of the present invention. [0229] The embodiment of FIG. 16 discloses a method in which prediction blocks within the same coding block share spatial merging candidates and temporal merging candidates when the size of the coding block is equal to or smaller than a specific size in inter-frame prediction using merge mode. [0230] Several pieces of information can be used to use a method of sharing a single merging candidate list only in blocks that satisfy a specific condition. For example, information on whether a current block uses a single merging candidate list or not can be derived based on information on the size of a block on which parallel merging processing can be performed and information on the size of a current coding block. Spatial merging candidates and temporal merging candidates for a prediction block can be derived on the basis of a coding block that satisfies the specific condition based on the pieces of derived information. [0231] Referring to FIG. 16(A) , only when conditions that the size of a block on which parallel merging processing can be performed is 8×8 or greater and the size of a coding block is 8×8 are satisfied, for example, prediction blocks partitioned from the coding block can share a single merging candidate list. [0232] It is assumed that a first coding block CU0 1600 has a size of 32×32, a second coding block CU1 1610 has a size of 16×16, a third coding block CU2 1620 has a size of 32×32, a fourth coding block CU3 1630 has a size of 16×16, and a fifth coding block CU4 1640 has a size of 8×8. [0233] FIG. 16(B) is a conceptual diagram only showing spatial merging candidate blocks for some coding blocks. [0234] Referring to FIG. 16(B) , the second coding block 1610 can be partitioned into two prediction blocks 1615 and 1618 having an nL×2N form, and the fifth coding block 1640 can be partitioned into two prediction blocks 1645 and 1650 having an N×2N form. In FIG. 16(B) , it is assumed that a single merging candidate list for only the coding block 1640 having the 8×8 size is generated. [0235] Each of the first prediction block 1615 and the second prediction block 1618 of the second coding block 1610 can derive spatial merging candidates for each prediction block and generate a merging candidate list for each prediction block based on the derived spatial merging candidates. [0236] The size of the fifth coding block 1640 is 8×8, and the fifth coding block 1640 can satisfy conditions of the size of a block on which parallel merging processing can be performed and conditions of the size of a current coding block. In this case, the third prediction block 1645 included in the fifth coding block 1640 and the fourth prediction block 1650 can generate a single merging candidate list based on the spatial merging candidates and the temporal merging candidates derived on the basis of the location and size of a coding block. Accordingly, the reference picture index of a temporal merging candidate can be derived as one value. [0237] The reference picture index of the temporal merging candidate can be derived by the aforementioned methods. [0238] The above-described image coding and image decoding methods can be implemented in the elements of the image coder and the image decoder described with reference to FIGS. 1 and 2 . [0239] Although the present invention has been described, a person having ordinary skill in the art will appreciate that the present invention may be modified and changed in various manners without departing from the spirit and scope of the present invention which are written in the claims below.

The present invention relates to a method and apparatus for setting a reference picture index of a temporal merging candidate. An inter-picture prediction method using a temporal merging candidate can include the steps of: determining a reference picture index for a current block; and inducing a temporal merging candidate block of the current block and calculating a temporal merging candidate from the temporal merging candidate block, wherein the reference picture index of the temporal merging candidate can be calculated regardless of whether a block other than the current block is decoded. Accordingly, a video processing speed can be increased and video processing complexity can be reduced.

FIELD OF THE INVENTION The present invention relates generally to an apparatus for tightening chains placed on vehicle tires and methods for using such an apparatus. In particular, the present invention relates to an apparatus for uniformly tightening the chains for motor vehicle tires and methods for using such an apparatus. BACKGROUND OF THE INVENTION During ice, snow or mud conditions, as well as other reduced traction conditions, chains are often installed on motor vehicle tires and trailers to increase traction. A problem often associated with the chains is that they may have a lot of slack and may be quite loose after they are installed. Examples of known devices for increasing the tension on such chains are described, for example, in U.S. Pat. Nos. 4,173,244, 4185,674, 4,237,951, 4,266,593, 4,392,521, 4,679,608, 4,799,522, 5,284,196, 5,785,783, 5,804,001, 6,026,876, and 6,085,816, the disclosures of which are incorporated herein by reference. Many of these chain-tightening devices are extremely complex and difficult to use. The non-uniformity of the tightening mechanism of other chain tighteners causes many problems. One such problem is chain roll. Chain roll occurs when loose portions of the chain bunch up and create their own rolling action. Such rolling action can lead to gouging and damaging of the tires, requiring the tires to be replaced. Also, the non-uniformity of the tightening mechanism can lead to chains that are installed too loosely. Many of the known chain tighteners are unable to keep the chain tight during operation, allowing the chain to work loose or shift off-center. In either instance, the loosened or off-center chain can fly up and damage to the vehicle. SUMMARY OF THE INVENTION The present invention includes an apparatus for tightening chains that are installed on a tire, as well as methods for using such an apparatus. The chain-tightening apparatus provides a mechanism for uniformly tightening the chains, as well as keeping a uniform tension on the chains while the tire rotates. The apparatus attaches to the chains quickly and easily, making the method of tightening such chains simple. The present invention comprises a system for tightening a tire chain or like antiskid devise upon a tire. The system comprises a plurality of attachments or hooks for attaching the apparatus to chain links or other suitable attachment points located at radially spaced intervals on the chain. A flexible connector or cable in the form of a closed loop is movably connected to the hooks. An apparatus for taking up slack, such as a reel with a winding inner cylinder, is used to take up slack in the cable and put the cable under tension. When the cable is tensioned there is an essentially constant tension around the loop of the cable imparting at each of the hooks a chain-tightening center-directed radial force. The reel can be locked to main maintain the cable in a tensioned state. When it is desired to remove the apparatus, the lock is released to release the tension on the cable. The attachment points are radially spaced, which means they are located around a central point, preferably symmetrically and essentially equally spaced. Prefect symmetry and spacing is not required, but spacing should be such that the chain is held in on the tire. It is an essentially symmetrical center-directed force that effectively holds the chain on the tire, and any asymmetry in the attachment of the apparatus of the invention should not significantly compromise the uniformity of the holding force on the chain. BRIEF DESCRIPTION OF THE DRAWINGS FIGS. 1-6 are views of chain-tightening apparatus and methods of using the same according to the present invention, in which: FIG. 1 is a schematic diagram showing one aspect of a chain tightening apparatus of the invention installed on a tire; FIG. 2 is schematic diagram of an apparatus of the invention when not mounted on a tire. FIG. 3 is a schematic showing the apparatus and tire as in FIG. 1 from a different view. FIG. 4 is a schematic of a slack take-up or tightening apparatus of the invention. FIG. 5 is a schematic similar to FIG. 2, showing another aspect of the invention using cable stops at the slack take-up apparatus. FIG. 6 is another schematic view of the slack take-up or tightening apparatus of the invention of FIG. 4 . FIG. 6A is a schematic view similar to FIG. 6, showing another aspect of the slack take-up apparatus. FIGS. 1-6A presented in conjunction with this description are views of only particular—rather than complete—portions of the chain tightening apparatus and methods of using the same. DETAILED DESCRIPTION OF THE INVENTION The following description provides specific details in order to provide a thorough understanding of the present invention. The skilled artisan, however, would understand that the present invention can be practiced without employing these specific details. Indeed, the present invention can be practiced by modifying the illustrated apparatus and method and can be used in conjunction with apparatus and techniques conventionally used in the industry. The chain tightening (CT) apparatus of the invention comprises at least three parts or components. The first part is an attachment for attaching the apparatus to the tire chain. The second part of the invention is a connector for connecting the attachment. The third part of the invention is a slack take-up for drawing up or tightening the connector. Reference is now made to FIG. 1, FIG. 2, and FIG. 3 . The chain tightening apparatus 5 of the invention comprises an attachment 20 , a connector 25 , and a slack take-up 30 . In use, as shown in FIG. 1, the chain tightening apparatus 5 is connected to the chains 10 on tire 15 . The apparatus of the invention can be employed on any type of chains 10 known in the art, such as chains for use in mud, ice and snow. Indeed, the apparatus of the invention can be used on any chains used to increase the traction of tires on any type of surface. The only requirement for the chains 10 is that, as described below, the attachment 20 of the chain tightening apparatus 5 is able to attach to the chains. Tire chains typically comprise two closed loops connected by cross-chains. When attached to the tire, the closed loops extend around the inside and outside sidewalls spaced from the tread with the cross-chains extending over the tread. The apparatus of the present invention is usually connected at spaced intervals on the loop on the outside sidewall. Accordingly the attachments 20 are spaced generally around a circumference of a circle defined by the loop. However, any chain construction that allows such radial spacing of the attachments is contemplated. Such construction includes “chains” or antiskid devices that are not constructed of linked chain structures, but can still provide the attachments as described herein and can be held and tightened on the tire by the tightening force of the invention. Any tire 15 known in the art can be employed in the invention. The apparatus of the invention 5 can be modified for any size and shape of tire by increasing (or decreasing) the number of attachments 20 and increasing (or decreasing) the length of the connection 25 described below. As well, the apparatus of the invention can be used with tires on self-propelled equipment and on motorized vehicles as well as on other wheeled apparatus pulled by motorized vehicles. For example, the chain tightening apparatus of the invention can be used with tires on cars, trucks, busses, semi-trucks and trailers, farm equipment, commercial equipment, lawn and garden tractors, and self-propelled snowblowers. To illustrate the invention, a tire with a width of about nine inches and a diameter of about 22.5 inches is depicted in FIG. 1 . As illustrated in FIG. 1, the CT apparatus 5 contains attachments 20 , connector 25 for connecting the attachments and slack take-up 30 . The function of an attachment 20 is to attach the CT apparatus to the chains 10 . The connector 25 connects the various attachments 20 together. The slack take-up 30 pulls on or tensions the connector 25 to draw up or tighten the connector. As mentioned above, the attachments 20 attach the CT apparatus 5 to the chains 10 . The attachments 20 can be attached to the desired portions of the chains 10 preferably after the chains are attached to the tire. The attachments 20 are removably attached to the chains so that the CT apparatus 5 can be used together with the chains or so that the chains can be used as-is without the CT apparatus. Once attached and the connector is tensioned, however, the attachments 20 should securely attach the CT apparatus 5 to the chains 10 . Any suitable attachments 20 known in the art accomplishing these functions can be employed in the invention. Examples of such attachments 20 include clips, hooks such as retractable hooks and s-hooks, latches and clasps. Preferably, s-hooks 20 are employed, as illustrated in the figures, as these securely attach to the chains while being removable. As depicted particularly in FIG. 3, the larger loop of an individual s-hook is connected to an individual link 40 of the chain 10 . When used in combination with the rest of the CT apparatus 5 , an attachment 20 pulls on link 40 which in turn pulls on other connecting links until the links are tight against tire 15 . The attachments 20 are connected to a plurality of locations on the chains 10 . The number of the attachments (and the corresponding number of attachment locations on the chains 10 ) depends on the circumference of the tire, the type (and size) of chains used, the type of attachment and the type of tire used. The number of attachments can typically range from 6 to 10 but a greater number of attachments could be used. In one aspect of the invention illustrated in FIG. 1, 7 s-hooks are employed in the CT apparatus of the invention. The attachments 20 are also connected to the connector 25 . While the attachments 20 and connector 25 can be removably connected, the attachments are preferably permanently connected in a secure fashion to the connector that allows the connector to easily slide through each attachment 20 . Any suitable connection system known in the art can be used between the attachment and the connector. The type of connection will depend in part on the attachment used in the CT apparatus. For example, when s-hooks are employed as the attachment, the smaller loop of the s-hook can enclose a cable-type connector as more fully described below. The connector 25 of the CT apparatus connects through all of the multiple attachments 20 . As described in more detail below, the connector 25 is movably connected to the attachments 20 . This is required so that when the connector 25 is tightened, the force from the tightening is more or less distributed evenly among the connector. While, an attachment may be fixed to the connector, preferably, no single portion of the connector is immovably attached to any given attachment 20 . An example is show in FIGS. 2 and 5, which show a cable connector slidably attached to s-hook attachments. Thus, the connector 25 exhibits the ability to move (or slide) with respect to any given attachment and corresponding chains 10 to which the attachments are connected. Any connector known in the art functioning in the above manner can be employed in the present invention. Examples of connectors that can be employed in the present invention include wires, straps, cording and cable. Preferably, a cable in the form of a loop is employed as the connector. As illustrated particularly in FIGS. 2 and 5 (showing s-hook attachments together), the cable-type connector 45 connects to all the s-hooks by passing through the s-hooks in a manner that allows the s-hooks to slide along the cable. The connector 25 is made of any suitable material known in the art that is sufficiently strong, e.g., not breakable when the CT apparatus is used in the manner described herein. Further, the material used for the connector should allow the attachments 20 to easily move along its length. It should be flexible so that it can have various configurations and can be transported easily. As well, the material for the connector should be lightweight. Any materials satisfying these criteria can be employed in the invention as the material for the connector. Exemplary materials for the connector include high-strength polymers, composite materials, metals, and twisted multiple-strand cable. ⅛ inch diameter galvanized aircraft cable employed as the material for the connector has been found suitable. The length of the connector is such that when attached through the attachments, and the connector is tightened, the combination functions to tighten the chains upon the tire. The size is generally determined by the projected size of the tires for which the chain tightening apparatus of the invention will be used. When cable is used as the connector, the length of the loop of the cable will typically be slightly smaller than the circumference of the tire. In one aspect of the invention, the circumference of the loop is about 5 feet to accommodate standard tire sizes for pick-ups, passenger cars, and sport utility vehicles. In another aspect of the invention, the circumference is about 8 feet to accommodate commercial vehicles, such as semi-trucks and trailers. The connector 25 of the CT apparatus 5 may be a permanently closed-end loop on which the slack take-up 30 is non-detachable by virtue of the connector running through a hole in the slack take-up, which in the illustrated slack take-up passes through the inner cylinder 55 . Such is the configuration for a standard 5 foot or 8 foot diameter permanently closed-end loop connector in the aspect of the invention illustrated in FIG. 2 . In another aspect of the invention, as illustrated in FIG. 5, the connector consists of an open-ended length of a suitable material for the connector 25 (such as a cable) which separately attaches to the slack take-up 30 using, for example, a slotted coupling device on one side of an outer cylinder walls 66 or through a slot 105 in the inner cylinder of the slack take-up 30 . The third component of the CT apparatus is the slack take-up 30 . by applying a rotational force to the slack take-up 30 , the slack take-up pulls on or tensions the connector 25 to draw up or tighten the connector and then retains the connector 25 in the tightened position while the CT apparatus remains installed on the chains and tire. Any device which accomplishes these functions can be employed as the slack take-up 30 in the invention. See, for example, the device described in U.S. Pat. Nos. 4,173,244, 4185,674, 4,237,951, 4,266,593, 4,392,521, 4,679,608, 4,799,522, 5,284,196, 5,785,783, 5,804,001, 6,026,876, and 6,085,816, the disclosures of which are incorporated herein by reference. In one aspect of the invention, the device depicted in the figures and shown particularly in FIG. 4 and FIG. 6 is employed as the slack take-up. This device comprises a take-up reel 50 containing an inner cylinder 55 enclosed by two outer circular walls 60 and 65 of a larger diameter. The inner cylinder contains a hole 70 running through its diameter through which a cable-type connector 25 passes. The inner cylinder can be any diameter but is preferably about 1¼ to 1½ inches in diameter. The diameter of hole 70 is sufficient to allow the cable loop to slide easily through the hole. The diameter of the outer cylinder walls can vary, but should be of sufficient diameter to take up and retain the slack in the connector. In the aspect of the invention illustrated in FIG. 4, the diameter of the outer cylinder walls is about 3 inches. An indentation 75 is located on the outer side of one of the outer cylinder walls at its center. This indentation is designed to operably connect with a tool, or the like, to rotate the slack take-up as described in more detail below. Thus, the indentation 75 can take on any number of shapes and sizes depending on the tool or mechanism used to rotate the slack take-up 30 . In one aspect of the invention, the indentation has a rectangular shape as depicted in FIG. 4 . Reference is now also made to FIG. 6 . An indented slot 105 extends on the outer side of one of the outer cylinder walls 65 to receive and retain the two ends of an open-ended connector 25 . This convention is used to allow any variable length of a connector 25 to be utilized by the CT apparatus in order to accommodate any non-standard tire sizes, such as tires on a large commercial farming tractor. Any closing structure for closing or securing the ends of an open loop cable to form a closed loop is contemplated by the present invention. In the aspect of the invention shown in FIG. 6, each of the open-ended cable ends of a connector are crimped with a stop 110 and cable ends with the stops are each inserted into the indented slot 105 to complete a secured loop of the connector. The connector slack is then drawn up on the reel 50 of the slack take up 30 . The stop 110 can be any suitable structure formed by crimping, molding, or bolting, or by any other way of attaching a stop to the cable end. In one aspect of the invention, the inner cylinder also contains a narrow slot 115 , extending along a major portion of the width of the inner cylinder through which the open ends of an open-ended connector consisting of a web or woven strapping material can be inserted in opposing directions to complete a secured loop of the connector. Reference is now made to FIG. 6 A. An end of the cable 45 is passed through the cylinder hole 70 and a stop 110 attached to the end to prevent it from passing back through the cylinder hole 70 . The other end of the cable is wrapped around the inner cylinder 55 of the take-up reel 50 a enough times to resist unwrapping of the cable (e.g., two or three times) and the free end passed through slot 105 . This secures the cable into a closed loop. The connector slack is then drawn up on the reel 50 of the slack take up 30 . Referring again to FIG. 6, the outer cylinder walls 60 and 65 contain a plurality of holes 80 around their outer circumference, with holes on one outer cylinder wall matching corresponding holes on the other cylinder wall to make sets of paired opposing holes 85 . Any number and location of holes 80 (and sets of paired holes 85 ) can be employed consistent with their purpose as described below. Located on opposing sides of any given outer cylinder wall are retaining pins 90 . The retaining pins have a straight side 95 that slides into any given set of paired holes 85 . The retaining pins 90 also contain an opposing side 100 that locks the retaining pin 90 into the set of holes 85 . Accordingly the pins are locked through the slack take-up 30 during a time the CT apparatus is installed on the tire, but can be manually inserted and removed by hand or by a tool when the tire is not rotating. The opposing side 100 can have any configuration—i.e., shape or size—accomplishing this locking function. In one aspect of the invention, the configuration of the opposing side is depicted in FIG. 6 . In this configuration, the opposing side has at least one narrow part in close proximity to the straight side 95 such that when inserted into the set of paired holes 85 , the narrow part abuts the edge of at least one outer cylinder. When forced past the edge, the retaining pin 90 thereby locks the take-up reel in that position until manually removed. The CT apparatus of the invention is used in the following manner. A tire is equipped with a chain as known in the art. A CT apparatus for the given size of the tire and the configuration of the chain is then attached by first attaching multiple attachments 20 on the connector 25 . For example, this is performed by providing the s-hooks on a connector 25 in the form of a cable 45 , which is shown in FIGS. 2, 4 , 5 and 6 . Next, the connector 25 is already attached to the slack take-up through the cylinder hole 70 if, for example, a closed-end connector 25 is used; or if an open-end connector securing the connector into a closed loop by any suitable method, such as by attaching the open ends to the indented slot 105 or passing the ends through the slot 115 in the inner cylinder in opposite directions. In each case, the cable is secured to form a completed loop. The assembled CT apparatus is then placed on chains 10 using the following procedure. The attachments 20 are first attached to the chain approximately equidistant around the circumference of the connector 25 . For example, the s-hooks can be attached to individual links of the chain at selected locations. The locations are selected for a uniform spacing between successive s-hooks. The slack take-up 30 is then used to draw up or tighten the connector 25 until the desired tension is reached. In the aspect illustrated in FIGS. 4 and 6, the take-up reel 50 is rotated, usually by means of a tool operably connected at indentation 75 , to tighten the connector until the attachments 20 pull the chains 10 into a tight position. Because of the sliding relationship of connector 25 to the attachments 20 , a uniform tension will result. When tightened in this manner, the take-up reel 50 is not at the center of the tire. Rather, the take-up reel is on the circumference of the connector similar to that depicted in FIG. 1 . The take-up reel 50 is then locked in the tightened position where uniform tension is distributed throughout the chain by reinserting the previously removed retaining pins 90 while holding the take-up reel with a tool that fits within indentation 75 to wind or draw the excess cable of the loop onto the inner cylinder of the take-up reel. The tool may be any suitable structure that provides the suitable leverage and fits in indentation, such as, for example, a ratchet. The reinserted retaining pins prevent the connector 25 from unwinding. When the tire rotates, the locked take-up reel retains the connector in its tightened position 30 , maintaining the uniform tension obtained when the take-up reel was tightened. When the user wishes to remove the tire chains, the CT apparatus is easily removed by removing the retaining pins, which allows the cable to loosen, permitting the s-hooks to be removed. The chain may then be removed from the tire. In an alternate aspect of the invention, the CT apparatus has a ratcheting mechanism built into the CT apparatus. In this instance removable pins would not be required. Instead, the mechanism would involve a take-up reel rotating on a second shaft that holds the ratcheting mechanism, with one end of a cable connector attached to an outer wall and the other end connected to the take up reel. The reel is rotated around the ratchet mechanism to take up the slack and would have a release device to allow the ratchet mechanism to release the cable so that the CT apparatus can be removed. Other suitable ratchet, winch, lever or other mechanisms that can take up slack and maintain the tension on the connector as described herein, are contemplated in the present invention. Unlike apparatus known in the prior-art, the CT apparatus of the invention achieves a uniform tension on the chains. The uniform tension is obtained since the connector is drawn linearly along its circumference creating a corresponding uniform pull of the tire chains in a radial direction through the attachments toward the center of the tire, like pulling the drawstring on a knapsack. The chain links not directly connected to attachments are likewise drawn tight via their linked interconnection to the attached links. None of the apparatus in the prior-art provide this same uniform radial pull by circumferentially tightening the connector. If desired, any metal parts comprising the CT apparatus can be plastic-coated to inhibit rust formation and other corrosion. The slack take-up can be made from high-impact plastic, metal, or any other suitable material. The size and strength of the individual components will depend upon the size of the tire and the weight of the chains that are installed on the tire. Having described the preferred embodiments of the present invention, it is understood that the invention defined by the appended claims is not to be limited by particular details set forth in the above description, as many apparent variations thereof are possible without departing from the spirit or scope thereof.

Apparatus for tightening chains that are installed on a tire, as well as methods for using such apparatus is disclosed. The chain-tightening apparatus provides a mechanism for uniformly tightening the chains, as well as maintaining a constant uniform tension on the chains when the tire moves. The apparatus attaches to the chains quickly and easily, making the method of tightening such chains simple. The apparatus is quickly and easily removed and is easily and compactly stored.

[0001] The following specification particularly describes the invention and manner in which it is to be performed: TECHNICAL FIELD OF INVENTION [0002] The invention discloses a reversible switching system that switches surfaces from being hydrophilic to hydrophobic and vice-versa. Particularly, the invention relates to reversible switching system comprising a positively charged conjugated polymer; a micelle and a hydrocarbon component that is affected by a potential such that the system switches from being super hydrophilic to super hydrophobic due to the polyvalent interaction between surfactant assemblies and delocalized charges of an insoluble polymer. More particularly, the present invention relates to a process for the switching of a system from being hydrophilic to hydrophobic and vice versa. These systems find application in antisticking coating, anticontamination coating and wicking surfaces. BACKGROUND AND PRIOR ART [0003] Polyelectrolytes have been widely used as scaffold for rendering mechanical stability to dynamic amphiphilic assemblies (AAs) and to prepare materials with specific architecture, nanoscopic containers and surfaces with low interfacial energy. The Coulombic attraction between oppositely charged polyelectrolytes cooperative aggregation process that eventually leads to materials with structural features similar to the AAs ( FIG. 1 a ). Furthermore, it has been found that the structural features of AAs remain unaltered upon interaction with biopolymer electrolytes having a semirigid backbone such as DNA. The common features of these polyelectrolytes are the charges, solubility in water, and coil-like structure in solution. In fact, the coil-like structure in solution permits these polyelectrolytes to adapt a new conformation and envelope the AAs ( FIG. 1 a ). [0004] There are couple of prior art relates to reversible switching for amphiphiles such as Article titled “Electrical Switching between Vesicles and Micelles via Redox-Responsive Self-Assembly of Amphiphilic Rod-Coils” by Hoon Kim in J. Am. Chem. Soc ., 2011, 133 (14), pp 5206-5209 discloses an aqueous vesicular system i.e. switchable by electric potential without addition of any chemical redox agents into the solution using redox-responsive self-assembly of an amphiphilic rod-coil molecule consisting of a tetraaniline and a poly(ethylene glycol) block, whereas optical switching of self-assembly and disassembly of noncovalently connected amphiphiles is disclosed in Langmuir 2007, 23, 12791-12794 by Jiong Zou et al. where a hydrophobic compound, i.e. 3C18-Azo, containing an azo head and three 18C alkyl chains to form an amphiphile by capping it with a cyclodextrin (CD) by inclusion complexation. Also Lahann J. et al. in Science 2003, 299, 371-374 reported the design of surfaces that exhibit dynamic changes in interfacial properties, such as wettability, in response to an electrical potential. [0005] Superhydrophobic surfaces with water contact angles above 150° are suitable for antisticking, anticontamination, and anticorrosion technologies. Similarly, superhydrophilic materials with water contact angles below 10° have applications as a wicking material in heat pipes. Ayaka Uyama in Langmuir , 2011, 27 (10), pp 6395-6400 discloses, reversible topographical changes on diarylethene microcrystalline surface between the rough crystalline surface of an open-ring isomer and flat eutectic surfaces by alternate UV and visible light irradiation wherein the contact angle changes of a water droplet between 80° and 150°. [0006] The wettability of a solid surface is strongly influenced both by its chemical composition and by its geometric structure (surface roughness). Several experiments have focused on exploiting surface roughness to engineer superhydrophobicity or superhydrophilicity. Microscale roughness features as well as nanoscale features have also been investigated. [0007] Reversible superhydrophobicity to superhydrophilicity transition by extending and unloading an elastic polyamide film is disclosed in Macromolecular Rapid Communications 26, (6), pp 477-480, 2005 by Jilin Zhang whereas reversible conversion of conducting polypyrrole (ppy) films from superhydrophobic to superhydrophilic” by Dr. Xu I. in Angewandte chemie international edition 44, (37), 6009-6012, 2005; additionally Stephen Feldberg in J. Am. Chem. Soc. 1984, 106, 4671-74 reported reinterpretation of polypyrrole electrochemistry consideration of capacitive currents in redox switching of conducting polymers. [0008] Further US2011244046 (Shen Youqing et al.) describes charge reversible polymers, peptides and their resulting colloidal particles, where charge reversal of the poly ([epsilon]-caproactone)-block-polyethyleneimine PCL-PEI/amide micelles was determined by measuring their [zeta]-potentials at different acidities. Further the zeta potential disclosed in state of art is more than +50V and −20V, that requires more energy and also, the reusability of the said polymers are not feasible. [0009] Further, the prior art fails to made out a simple approach that would create both chemical and structural features to show superhydrophilic and superhydrophobic nature of the surface. A surface that can switch between hydrophilic to hydrophobic will have a profound impact on preparing smart surfaces. The approach adopted so far is tedious. Till date, the work concerning such superhydrophilic to superhydrophobic switching has used patterned surfaces to increase the roughness in a calculated manner. [0010] In view of the foregoing, there is still need in the art to provide an simple and improved system that can switch between hydrophilic to hydrophobic at lower potential, useful in antisticking coating, anticontamination coating and wicking surfaces. OBJECTIVE OF THE INVENTION [0011] The main objective of the invention is to provide a reversible switching system that can switch from being hydrophobic to hydrophilic and vice versa with a contact angles of <10° and >150°. [0012] Another objective of the invention is to provide a process for the switching of a system from being hydrophobic to hydrophilic and vice versa. SUMMARY OF THE INVENTION [0013] The present invention provides a reversible switching system comprising a positively charged insoluble conjugated polymer; a micelle and a hydrocarbon component that is affected by a potential such that the system switches from being super hydrophilic to super hydrophobic and vice versa, due to the polyvalent interaction between surfactant assemblies and delocalized charges of the insoluble polymer. [0014] In one embodiment of the present invention the insoluble polymer is oxidized conjugate polymer. [0015] In another embodiment of the present invention the oxidized conjugate polymer is polyaniline. [0016] In another embodiment of the present invention the micelle is SDS micelle. [0017] In another embodiment of the present invention the SDS micelle is optionally encapsulated with guest molecule selected from chemical, dyes. [0018] In another embodiment of the present invention the hydrocarbon component is selected from naphthalene or behra's amine. [0019] In still another embodiment of the present invention the behra's amine is di-tert-butyl 4-amino-4-(3-tert-butoxy-3-oxopropyl)heptanedioate (branched amino alkyl ester). [0020] In yet another embodiment of the present invention A process for the switching of a system from being hydrophilic to hydrophobic and vice versa, which comprises: a) Providing a positively charged insoluble conjugated polymer followed by dipping the same in micelle solution to obtain a hydrophobic surface; b) Filling the hydrophilic gaps between two micelle chains with a hydrocarbon component to obtain superhydrophobic surface; and c) Neutralizing the surface by applying a reduction potential at a voltage range of −0.2 to eject SDS, the aromatic and aliphatic hydrocarbons and thereby to obtain the superhydrophilic surface. [0024] In yet another embodiment of the present invention the provision of positively charged conjugated polymer comprises the oxidation of the polymer at a voltage range of 0.2 to 0.6V. [0025] In yet another embodiment of the present invention the filling of the hydrophilic gaps between two micelle chains comprises dipping the hydrophobic polymer surface in branched alkyl chain hydrocarbons. [0026] In yet another embodiment of the present invention the insoluble polymer is oxidized conjugate polymer. [0027] In yet another embodiment of the present invention the oxidized conjugate polymer is polyaniline. [0028] In yet another embodiment of the present invention the micelle is SDS micelle optionally encapsulated with guest molecule selected from chemical, dyes etc. [0029] In yet another embodiment of the present invention the hydrocarbon component is selected from naphthalene or behra's amine. [0030] In yet another embodiment of the present invention the behra's amine is di-tert-butyl 4-amino-4-(3-tert-butoxy-3-oxopropyl)heptanedioate (branched amino alkyl ester). DESCRIPTION OF DRAWINGS [0031] FIG. 1 depicts a) polyelectrolyte enveloping AAs. b) rigid OCPs disassembling AAs. c) The size of SDS AAs at critical micelle concentration as determined by DLS. d) Unaltered size of SDS AAs upon exposure to NCPs as determined by DLS. [0032] FIG. 2 depicts a) Comparison of percentage of pyrene released upon interaction with OCPs bulk powder and OCPs nanofibers. b) TEM image showing the morphology of OCPs nanofibers. c) BET curve used to calculate the surface area of nano OCPs and d) bulk OCPs. [0033] FIG. 3 depicts a) reversible assembly and disassembly of AAs as a function of applied potential. b) Fluoresence emission spectra of pyrene in an aqueous phase (I372/I384=1.45) upon disassembly by OCPs film. c) Fluoresence emission spectra of pyrene encapsulated in SDS AAs that are formed from the surfactants ejected from the OCPs surface (I372/I384=0.94). [0034] FIG. 4 depicts a) Percentage of pyrene released from SDS AAs upon disassembly induced by OCPs film. b) Percentage release of pyrene as a function of the number of surface regenerations. [0035] FIG. 5 (I) depicts a) Digital image showing water droplet on OCPs surface. b) Image showing water droplet on NCPs surface. c) CA of water droplet on CPs oxidized at 0.2 V and treated with SDS AAs. d) CA of water droplet on CPs oxidized at 0.4 V and treated with SDS AAs. e) CA of water droplet on CPs oxidized at 0.6 V and treated with SDS AAs. f) CA of water droplet on CPs oxidized at 0.6 V, treated with SDS AAs, and subsequently treated with Behera's amine. [0036] FIG. 5 (II) depicts change in water contact angle (°) as a function of applied potential (V), sodium dodecyl sulfate (SDS) and aromatic and aliphatic hydrocarbon. The change in the contact angle with the addition of a micelle (SDS) and a hydrocarbon is represented. [0037] FIG. 6 depicts a) AFM image showing the morphology of OCPs surface b) AFM image of the surface of OCPs treated with SDS AAs. c) AFM image of OCPs treated with SDS AAs, which was then treated with Behera's amine. [0038] FIG. 7 depicts superhydrophilic surface. [0039] FIG. 8 depicts hydrophobic surface. ABBREVIATIONS [0040] AAs: Amphiphilic Assembly (herein SDS micelle assembly) CPs: Conjugated polymers OCPs: Oxidized conjugated polymers NCPs: Neutral conjugated polymers CA: Contact angle DLS: dynamic light scattering SDS: sodium dodecyl sulfate CTAB: cetryltrimethylammonium bromide CMC: critical micelle concentration AFM: Atomic force microscopy DETAILED DESCRIPTION OF INVENTION [0041] When water contact angle is >150° and up to 180°, e.g. as in silane material, it is termed as super hydrophobic materials, the contact angel <10° is termed as super hydrophilic materials. [0042] Therefore, the need is for providing a surface with switching between 10 deg and 150 deg ie between superhydrophilic and superhydrophobic for industrial use. [0043] In a preferred embodiment, the invention provides a reversible switching system that comprises a positively charged conjugated polymer; a micelle (SDS) and a hydrocarbon component that is affected by a potential such that the system switches from being super hydrophilic to super hydrophobic and vice versa, due to the polyvalent interaction between surfactant assemblies and delocalized charges of an insoluble polymer. [0044] In another preferred embodiment, the invention provides a process for the switching of a system from being hydrophilic to hydrophobic and vice versa, which comprises: a) Providing a positively charged conjugated polymer followed by dipping the same in micelle solution to obtain a hydrophobic surface; b) Filling the hydrophilic gaps between two micelle chains with a hydrocarbon component to obtain superhydrophobic surface; and c) Neutralizing the surface by applying a reduction potential so as to eject SDS, the aromatic and aliphatic hydrocarbons and thereby to obtain the hydrophilic surface. [0048] Providing positively charged conjugated polymer comprises the oxidation of the polymer at a voltage range of 0.2 to 0.6V. [0049] Filling the hydrophilic gaps between two micelle chains comprises dipping the hydrophobic polymer surface in branched alkyl chain hydrocarbons. [0050] To accomplish the objective of the invention, a charged polymer was taken and placed a water drop on it, which sticks at an angle of 6° ( FIG. 7 ). The polymer is dipped in a micelle, preferably SDS to obtain a contact angle of 100° ( FIG. 8 ). The increase is because the anionic head group of the SDS sticks to the positively charged surface. This increase in contact angle indicate the conversion of the hydrophilic surface to hydrophobic. However, the inventors have observed certain gaps, between two SDS chains which are hydrophilic. In order to make this surface super hydrophobic, the gaps must be filled with hydrophobic components. [0051] The process of the invention accomplishes with the filling of gaps between two SDS chains with a hydrocarbon, so as to increase the surface hydrophobicity up to 120°. The inventors surprisingly found that by filling the gaps with aliphatic branched hydrophobic molecule the contact angle can be further increased up to 155°. Then the polymer can be neutralized by applying a reduction potential so as to eject SDS, the aromatic and aliphatic hydrocarbons and to obtain the hydrophilic surface. This process can be repeated; hence the surface hydrophobicity and hydrophilicity can reversibly be obtained. [0052] With reference to FIG. 5 , as synthesized polymer exposed to water drop would show a contact angle of 6°. A polymer oxidized at 0.2 V and dipped in SDS micelles show 30°. Polymer oxidized at 0.4 V and dipped in SDS would show 60°. Polymer oxidized at 0.5 V and dipped in SDS would show 80°. Polymer oxidized at 0.6 V and dipped in SDS would show 100°. Polymer oxidized at 0.6 V and dipped in SDS and then dipped in aromatic hydrocarbon naphthalene shows 140°. Polymer oxidized at 0.6 V and dipped in SDS and then dipped in branched alkyl chain hydrocarbon shows 155°. [0053] Thus in accordance with the objectives of the invention, the invention provides electropolymerisation of a monomer on its surface and keep the polymer in it's charged state. This surface is hydrophilic due to the presence of positive charges. Then the electrode is dipped in a solution of negatively charged amphiphilic assemblies, which get adsorbed on the polymer surface and convert the polymer surface to hydrophobic. [0054] In a preferred embodiment of the invention, the polymer is insoluble oxidized conjugated polymer, more preferably polyaniline. The micelle is SDS and the hydrocarbon component is naphthalene or behra's amine, wherein behra's amine is di-tert-butyl 4-amino-4-(3-tert-butoxy-3-oxopropyl)heptanedioate (branched amino alkyl ester). [0055] To accomplish reversible assembly and disassembly of AAs, the polymer needs to be coated on substrate and charge regeneration has to be carried out. Accordingly, polyaniline was electrodeposited on Pt foil by applying a constant potential of 0.6 V vs Ag/AgCl. At this potential, the polymer bears delocalized positive charges; hence it is a film version of OCPs. The polymer film was then dipped in a solution containing SDS AAs, and the size of the assemblies in an aliquot was monitored by DLS. The disassembly of AAs was confirmed by the disappearance of a peak corresponding to 6 nm, and then a reduction potential of −0.2 V was applied with respect to a quasi-reversible Pt reference electrode. This potential is suffice to convert OCPs to NCPs that in turn would result in the ejection of SDS amphiphiles into the solution from the polymer surface. Now, the concentration of SDS in the solution would reach the critical micelle concentration (CMC), if all the SDS molecules that were bound to the OCPs surface were released. To check the formation of the AAs, aliquot withdrawn from the solution was subjected to DLS analysis and noted that a peak corresponding to 6 nm in the DLS histogram indicating the formation of SDS AAs. The cartoon in FIG. 3 a depicts the whole process of assembly and disassembly as a function of applied potential. [0056] In yet another embodiment, the invention provides the release and encapsulation of guest molecules upon disassembly and assembly of SDS AAs. Accordingly, OCPs film coated electrode was immersed in pyrene encapsulated SDS AAs solution and left quiescent. Emission spectrum of the solution was recorded to determine the environment of pyrene. The I 372 /I 384 was found to be 1.45, which indicates the disassembly of SDS AAs and the presence of pyrene in an aqueous environment ( FIG. 3 b ), and then the OCPs film was reduced to NCPs, which ejects the SDS surfactants into the solution. Upon reaching the CMC, the SDS AAs are likely to sequester the pyrene in the solution. If that occurs, the I 372 /I 384 of pyrene should indicate the presence of the guest molecule in a hydrophobic environment. The intensity of I 372 /I 384 of 0.94, which confirms the presence of pyrene in a hydrophobic environment ( FIG. 3 c ). [0057] With these experiments, the inventors have demonstrated the reversible assembly and disassembly of AAs and corresponding encapsulation and release of guest molecules. [0058] The release of guest molecules by reusing the OCPs surface and the release of pyrene from AAs are also demonstrated in the present invention. The change in I 372 /I 384 confirms the disassembly and concomitant release of the guest molecules from the AAs upon interaction with OCPs coated substrate. UV-vis absorption spectroscopy was used to quantify the released pyrene, and it was found to be 93% ( FIG. 4 a ). The polymer coated Pt foil was subsequently reduced at −0.2 V to convert the polymer to its neutral state (NCPs), which is not capable of inducing disassembly. Then the polymer film was reoxidized at 0.6 V and utilized for further disassembly of AAs. The experiment and the analysis were repeated to elucidate the films reusability. Thus the invention is made it possible to reuse the electrode over a period of five days and over ten cycles for the purpose of disassembling and releasing the guest molecules, after which the quantity of the release decreased, possibly owing to the degradation of the polymer ( FIG. 4 b ). The CTAB AAs used for control experiments did not disassemble, which further confirms the methods specificity. [0059] AAs assembly and disassembly, release of payload, and reuse of OCPs surface for repetitive disassembly have been unambiguously proven by the present invention. [0060] The present invention further demonstrates the change in surface properties of OCPs upon interaction with AAs. The OCPs surface is likely to be hydrophilic due to the presence of positive charges on the polymer. It is necessary to recall that the OCPs surface was prepared by oxidizing the polymer at 0.6 V vs Ag/AgCl, wherein the polymer is likely to have a maximum number of positive charges. On the other hand, NCPs have no positive charges because they are prepared by applying a reduction potential of −0.2 V vs Ag/AgCl to OCPs. One of the attractive features of CPs is it provides a handle to control the number of positive charges by varying the applied potential. Thus, by controlling the positive charges on the polymer, one can control the number of anions that bind with the polymer surface. The number of anionic surfactants that bind with the polymer is a function of the number of positive charges generated by the applied potential. Considering this, it is reasonable to anticipate that the hydrophobicity of the surface scales as a function of the amount of surfactant binding on the polymer surface. To prove this, the inventors have studied the contact angle (CA) of a drop of water on the polymer film surface, which gives a measure of surface hydrophobicity. The CA on the OCPs surface was found to be 6°, a value typical of hydrophilic surfaces ( FIG. 5 a ). Also, the CA of NCPs, which was found to be about 6° ( FIG. 5 b ), and then the oxidation potential of 0.2 V was applied to NCPs which generates positive charges on the NCPs surface. After that, the film was immersed in SDS AAs, and then the CA of water drop was measured. The CA for this surface was found to be 30° ( FIG. 5 c ). By repeating the same procedure, but by applying a potential of 0.4 V, a surface with CA of 60° was obtained ( FIG. 5 d ). This was further increased to 80° by applying 0.5 V to NCPs and dipping in SDS AAs. By applying 0.6 V and treating the film with SDS AAs a contact angle of 100° was achieved ( FIG. 5 e ). Upon interaction of SDS AAs with OCPs surface, the positive charges are neutralized by the negative charges of the SDS surfactants and concurrently the alkyl chains protrude from the polymer surface. This leads to low surface energy and results in an increase in hydrophobicity. It is enticing to increase the CA to ≧150° (superhydrophobic surface). OCPs surfaces prepared by applying potentials above 0.6 V and dipping in SDS AAs did not increase the CA beyond 100°, due to the formation of fully oxidized polyaniline (OCPs) at 0.6 V. Surprisingly, the inventors have found that the CA can be increased by filling the hydrophilic voids between the SDS molecules. This was achieved by immersing a OCPs film with CA of 100° in a solution containing di-tert-butyl 4-amino-4-(3-tert-butoxy-3-oxopropyl)heptanedioate (branched amino alkyl ester or Behera's amine. The branched structure was chosen to impart hydrophobicity to the hydrophilic voids between the SDS molecules. The CA increased from 100° to 155° ( FIG. 5 f ). The OCPs treated with branched amino alkyl ester alone showed a CA of 50°, which is 105° less than that observed for OCPs treated with SDS and Behera's amine. Furthermore, the CA was found to be 100°, if the acid analogue of Behera's amine was used. The carboxylic acid and amine terminals cannot impart hydrophobicity to the voids, hence the CA remain close to OCPs surface treated with SDS AAs. These experiments confirm that the filling of the hydrophilic voids between the alkyl chains of SDS molecules protruding from the OCPs surface as a means to increase hydrophobicity. Often, hydrophobicity enhancement has been attributed to increase in the surface roughness, hence it is necessary to test any microscopic changes on the polymer surface upon interaction with SDS AAs. Atomic force microscopy (AFM) imaging of OCPs ( FIG. 6 a ), OCPs treated with SDS AAs ( FIG. 6 b ), and OCPs treated with SDS AAs and Behera's amine ( FIG. 6 c ) was carried out to find out any changes in the morphology of the surfaces. No discernible changes in the morphology were observed in AFM images confirming the change in surface hydrophobicity is due to the noncovalent attachment of SDS surfactants and Behera's amine. [0061] Till date, CPs based superhydrophobic surfaces have been prepared using tedious synthesis and patterning procedures. However, the invention demonstrates that the surface wettability of CPs can be controlled at ease by varying the applied potential and dipping in SDS AAs. [0062] In a nut shell, the present invention utilized a charge bearing, insoluble, and rigid conjugated polymer to disassemble AAs. Furthermore, the disassembly and assembly of AAs were accomplished by switching CPs between its charged and neutral states. The rate of disassembly was controlled by modulating the morphology of the CPs. During the disassembly of AAs, the encapsulated cargo can be released, and the process can be repeated several times by regenerating the charges on the CPs. Upon disassembly, the amphiphiles which were constituents of AAs are bound to the OCPs that impart hydrophobicity to the polymer. By controlling the charge on the polymer and subsequently dipping in AAs, the surface hydrophobicity was varied between hydrophilic to hydrophobic wettings. [0063] The following examples, which include preferred embodiments, will serve to illustrate the practice of this invention, it being understood that the particulars shown are by way of example and for purpose of illustrative discussion of preferred embodiments of the invention. EXAMPLES Materials and Methods [0064] The analytical grade chemicals aniline, ammonium persulfate, sodium dodecyl sulfate (ACS reagent), cetyl trimethyl ammonium bromide, and pyrene were purchased from Sigma-Aldrich and used as received. Reagent grade HCl and isopropyl alcohol were purchased from Loba Chemie. The deionized water was collected from Millipore Q Gard water purifier and further purified by filtering through a 30 nm polycarbonate membrane purchased from SPI pore. The 30 nm pore membrane was mounted on a stainless steel filter holder, which was then fitted in plastic syringe. A CH Instruments 600D potentiostat/galvanostat was used for electrochemical measurements and electropolymerization of aniline. The working and counter electrodes were fabricated using Pt foil (99.9% purity) purchased from Arora Matthey Ltd. UV-vis spectra were recorded with a Jasco U Best V-570 UV-vis spectrophotometer. Fluorescence spectra were recorded with a Cary Eclipse Fluorescence spectrophotometer. The particle size of the assemblies was recorded with a Brookhaven 90 plus Particle Size Analyzer. AFM images were recorded with MM AFM LN supplied by Veeco Multimode in taping mode. Gold coated silicon substrates were used as substrates for AFM imaging. TEM imaging was done with a Jeol 1200 EX transmission electron microscope. The carbon coated copper grids (400 grids) were obtained from Ted Pella. Water drop contact angles were measured in a Digidrop Contact Angle Meter. Example 1 [0065] Polyaniline of ˜1 micron thickness was chosen as the polymer surface. This was synthesized on top of a gold coated plastic substrate. The water contact angle was found to be 6°. To modify the surface, SDS micelles were chosen (concentration 10 −3 M). The polymer was oxidized at various potentials and dipped in 10 −3 M SDS. This gave rise to change in surface hydrophobicity. To further enhance the surface hydrophobicity, the polymer surface was treated with 0.01 M behra's amine. Example 2 Synthesis of Polymer [0066] Oxidized conducting polymer (OCPs), polyaniline, was synthesized by chemical polymerization of aniline using ammonium persulfate (APS) as initiator. For this purpose 100 mM aniline was dissolved in 1 M HCl solution (100 ml), and subsequently 100 mM APS (ammonium persulfate) was added. The reaction was allowed to proceed for 24 h at 25° C., and then the reaction was stopped by filtering the polymer using 200 nm pore nylon membrane, followed by washing the initiator and unreacted monomer. Polyaniline nanofibers were synthesized by following the interfacial polymerization procedure. Alternately, aniline (100 mM) was dissolved in chloroform (100 ml) and allowed to stand quiescent in a sample vial. To that solution, 100 mM APS dissolved in 1 M HCl (100 ml) was added. After five minutes, the polyaniline formation at the chloroform water interface is visible. The reaction was allowed to proceed for 24 h at 25° C., and then the fibers were collected by filtration in a 200 nm pore nylon membrane. The polymer nanofibers were then washed with copious amounts of water to remove unreacted monomer and initiator. [0067] To prepare NCPs (Neutral conjugated polymers), 200 mg of OCPs was stirred in 25 mL of hydrazine hydrate for 12 h at 25° C. Hydrazine hydrate reduces the OCPs to NCPs. The NCPs were then washed thoroughly with water. Similarly, NCPs nanofibers were prepared by treating the OCPs nanofibers with hydrazine hydrate. Example 3 Preparation of Pyrene Encapsulated Micelles [0068] The micelles were prepared by dissolving 6 mM SDS (Sodium dodecyl sulfate), in deionized water (100 ml), which was prefiltered through 30 nm polycarbonate membrane. The pyrene encapsulation was carried out by following the reported procedure, and the concentration of pyrene was maintained at 10 −6 M for all the experiments. Pyrene exhibits multiple emission peaks, and the ratio of the peak intensity at 372 nm (I372) to that at 384 nm (I384) provides information about the environment of the probe. For this study, pyrene was entrapped in SDS AAs, and I372/I384 was determined to be 0.97, which corresponds to the presence of pyrene in a hydrophobic environment. [0069] To study the disassembly and release of guest molecules, 100 mg of OCPs was added to the SDS AAs solution, and the mixture was allowed to stand quiescent. For the DLS measurement, the supernatant liquid was taken out using a syringe and filtered through a 600 nm polycarbonate membrane (SPI pore). The size of the assemblies is of 6 nm, which is 100 times smaller than the pores (600 nm) of the membrane the size of the assemblies is not affected by this filtration step as evidenced from the DLS data. For UV-vis spectra and fluorescence emission spectra, the supernatant solution was taken out without disturbing the underlying insoluble polymer (either OCPs or NCPs). The TEM images of the polyaniline were obtained by dropping the sample, which was dispersed in isopropyl alcohol on a carbon coated copper grid (400 grid). OCPs/polyaniline films were prepared by electropolymerization of 0.1 M aniline dissolved in 0.1 M HCl by applying a constant potential of 0.6 V vs Ag/AgCl reference electrode. The working (geometric area—1 cm2) and counter electrodes (geometric area—2 cm2) were Pt foils. [0070] For reversible disassembly and assembly studies, the OCPs film coated electrode was dipped in pyrene encapsulated SDS AAs. After 3 h, an aliquot was withdrawn and the emission spectrum was recorded. The solution was poured into the solution with the OCPs coated electrode. A potential of −0.2 V was applied to eject the SDS unimers from the OCPs surface into the solution. The aliquot was withdrawn to record the emission spectrum to show the formation of SDS AAs ejected from the OCPs that eventually encapsulated the pyrene in the water. Example 4 Disassembly of AAs [0071] Sodium dodecyl sulfate (SDS) is chosen as the amphiphile because of the presence of anionic moieties which can, therefore, interact with the positive charge bearing OCPs. [0072] The OCPs induced disassembly process was first studied using dynamic light scattering (DLS). The aggregate size of AAs prepared using 6 mM SDS solution was found to be 6 nm ( FIG. 1 c ). In order to monitor the disassembly of AAs, 100 mg of dry oxidized polyaniline (OCPs) powder was added to 5 mL of SDS solution. Upon addition of OCPs, the peak corresponding to 6 nm disappeared, which indicates the disassembly of AAs. If the disassembly is due to the attractive interaction between the positive charges of OCPs and negative charges of the SDS AAs, the assemblies should not disaggregate in the presence of neutral conducting polymers (NCPs). To examine this possibility, NCPs was added to the AAs solution, and the assembly size was determined by DLS. The size of the assemblies was found to be 6 nm and remained constant over a period of 48 h ( FIG. 1 d ). This confirms that the disassembly process is exclusively due to the electrostatic attractive interaction between OCPs and SDS AAs. [0073] To further validate the above observation, cationic micelle, cetyltrimethylammonium bromide (CTAB) was allowed to interact with OCPs, and the assembly size was monitored. In this experiment, both the AAs and the OCPs have positive charges, hence, the electrostatic attraction is unlikely. Indeed, the DLS measurement indicated that the 6 nm CTAB assemblies did not undergo any alteration. Similarly, the CTAB micelle size was unaffected while it was exposed to NCPs, which indicates that the assembly process is solely due to the polyvalent attractive interaction between negatively charged AAs and positively charged rigid polymer Example 5 Pyrene Encapsulated AAs [0074] The disassembly and concomitant release of sequestered guest molecules were determined as a function of time using pyrene, which exhibits multiple emission peaks, and the ratio of the peak intensity at 372 nm (I372) to that at 384 nm (I384) provides information about the environment of the probe. [0075] For this study, pyrene was entrapped in SDS AAs, and I372/I384 was determined to be 0.97, which corresponds to the presence of pyrene in a hydrophobic environment. Upon addition of OCPs to the pyrene sequestered SDS AAs the I372/I384 started increasing from 0.97 and reached a value of 1.5 in 60 min, which indicated that the new environment faced by pyrene was hydrophilic. This change in the environment of pyrene is due to the disassembly of SDS AAs and concomitant release of cloistered pyrene into water. [0076] The same experiment was carried out using NCPs, and no alteration in I372/I384 of pyrene was observed, which corroborates the fact that the assemblies remain intact in the presence of NCPs. [0077] UV-vis spectroscopy was used to quantify the amount of pyrene released due to the disassembly of the SDS AAs upon interaction with OCPs. The gradual decrease in the absorption intensity at 338 nm indicates the release of pyrene from SDS AAs into water. About 96% of pyrene was released upon the disassembly of AAs, and the release profile is shown in FIG. 2 a . On the other hand, the unaltered UV-vis spectral features of pyrene indicate that the SDS AAs did not disassemble, and the guest molecules were not released in the presence of NCPs. From these experiments, it is apparent that the AAs disassemble only when involved in electrostatic attraction with insoluble and rigid OCPs surface. The release percentage value is also indicative of the fact that the released pyrene is essentially in the aqueous environment and not entrapped in the amphiphile modified CPs. This further proves the absence of any reformed assemblies comprising OCPs and amphiphiles. If it is a surface bound phenomenon, the release kinetics should vary with changes in the OCPs surface morphology. In order to verify this, OCPs nanofibers were synthesized by interfacial polymerization, having a diameter of about 50 nm ( FIG. 2 b ). [0078] Contrary to this, polyaniline synthesized by conventional solution polymerization yielded large quantities of bulk powder. The nanofibers synthesized by interfacial polymerization were then added to the SDS AAs, and the dimension of the assemblies was determined. The 6 nm SDS AAs were disassembled as can be seen from the DLS correlation data. To study the rate of the release, these nanoscopic fibers were added to pyrene encapsulated SDS AAs solution and release of pyrene was monitored as a function of time. [0079] The disassembly and concomitant release was completed in 15 min, which is faster by about three times as compared to the bulk OCPs release ( FIG. 2 a ). Usually, change in morphology from bulk to nano would increase the surface area significantly. [0080] Nanofibers induced guest molecules release should have been much faster than that observed in this study. To further investigate this, surface area of bulk OCPs and nano OCPs were determined by BET surface area measurements and found that the surface area of nano OCPs and bulk OCPs was found to be 45 m2/g ( FIG. 2 c ) and 20 m2/g ( FIG. 2 d ), respectively, which suggests the possibility of the presence of nanofibers in bulk OCPs. Indeed, meticulous TEM imaging of bulk OCPs showed the presence of nanofibers. It has also been reported that the polyaniline synthesized by conventional methods did contain nanofibers along with bulk powder. Consolidating the results of surface area measurements and guest molecule release studies, the ˜65% decrease in release time is a result of ˜50% increase in surface area of the nanofibers. [0081] From the foregoing, it is evident that the disassembly and release of payload occurs only when electrostatic attraction between rigid OCPs and AAs is in operation. Thus, the mere presence of CPs is not adequate to trigger disassembly and release of guest molecules. This essentially means that the CPs can be placed at a location passively and activated by generating positive charges on the CPs whenever desired. ADVANTAGES OF THE PRESENT INVENTION [0082] The present invention is very simple process by applying a small potential to change wettability from super hydrophilic top super hydrophobic in between.

The present invention disclosed herein is a reversible switching system that switches surfaces from being hydrophilic to hydrophobic and vice-versa. More particularly, the invention relates to reversible switching system comprising a positively charged conjugated polymer; a micelle and a hydrocarbon component that is affected by a potential such that the system switches from being super hydrophilic to super hydrophobic due to, the polyvalent interaction between surfactant assemblies and delocalized charges of an insoluble polymer results in surface modification of the polymer. These systems find application in antisticking coating, anticontamination coating and wicking surfaces.

CROSS REFERENCE TO RELATED APPLICATIONS This application is related to: U.S. patent application Ser. No. 08/673,254 entitled, "Dark field, photon tunneling imaging systems and methods," U.S. patent application Ser. No. 08/673,252 entitled, "Dark field, photon tunneling imaging systems and methods for measuring flying height of read/write heads," and U.S. patent application Ser. No. 08/671,708 entitled, "Dark field, photon tunneling imaging systems and methods for optical recording and retrieval," all by John M. Guerra. BACKGROUND OF THE INVENTION 1. Field Of The Invention This invention in general relates to the fields of imaging and metrology of surfaces and, more specifically, to systems and methods that use evanescent field illumination in real-time, whole-field imaging and measurement of surfaces at low magnification and at wide angle. 2. Description Of The Prior Art Photon tunneling microscopy for measuring and visualizing submicron surface topographic features is known. Descriptions of the use of photon tunneling microscopy are contained, for example, in: Hatrick, N. J., "Use of Frustrated Total Internal Reflection to Measure Film Thickness and Surface Reliefs," J. Appl. Phys., 1962.33: p. 321; McCutchen, C. W., "Optical Systems for Observing Surface Topography by Frustrated Total Internal Reflection and by Interference," The Review of Scientific Instruments, Vol. 35, p. 1340-45, 1964; Guerra, J. M., "Photon tunneling microscopy," in Proceedings from Surface Measurement and Characterization Meeting, Hamburg, SPIE Vol. 1009, September 1988, pp. 254-62, U.S. Pat. No. 4,681,451 entitled, "Optical proximity imaging method and apparatus," issued to Guerra, J. M. and Hummer, W. T. Jul. 21, 1987; U.S. Pat. No. 5,349,443 entitled, "Flexible transducers for photon tunneling microscopes and methods for making and using same," issued to Guerra, J. M. Sept. 20, 1994; and U.S. Pat. No. 5,442,443 entitled, "Stereoscopic Photon Tunneling Microscope," issued to Guerra, J. M. Aug. 15, 1995. Hatrick, McCutchen, Guerra/Plummer, and Guerra disclose whole-field reflected evanescent light microscopes where the sample is not transilluminated nor scanned, but rather is illuminated by an evanescent field from an unrestricted total reflection surface at the object plane of an epi, or reflected light, illuminator. Here, the sample can be opaque or transparent, thick or thin, and can be viewed in real-time with high energy throughput. Such microscopes are very sensitive to smooth surfaces because of their use of the exponentially-varying amplitude of the evanescent field in the vertical direction to sense very small surface height variation. On the other hand, rougher surfaces scatter light back into the microscope, which decreases contrast and sensitivity. Also, the deeper topography is rendered as bright, because these areas penetrate the evanescent field to a small degree so that the epi-illumination is nearly totally reflected. The difficulty in detecting and measuring small changes in bright scenes limits the observable topographic depths to about 3/4 of the illuminating wavelength (which is the wavelength in air divided by the index of refraction, n, and the sine of the angle of incidence, θ). Further, the illumination and imaging optics are coupled because the objective serves as the condenser as well. This limits in a practical sense the use of such instruments to the availability of suitable commercial objectives, magnifications, fields of view, and numerical aperture. In addition, it is difficult, because of the coupling of imaging and illumination optics, to affect the polarization, phase, incident angle, and direction of the illumination. This, in turn, restricts the ability to maximize the tunneling range, increase lateral resolution, or tunnel through less rare media such as water in biological applications. Devices in which evanescent light from transilluminated samples is scattered into objective pupils are described in G. J. Stoney, "Microscopic Vision," Phil. Mag. p. 332, at 348-49, 1896; Surface contact microscope, Taylor & Francis; Ambrose, E. J., "A Surface Contact Microscope for the Study of Cell Movements," Nature, Nov. 24, 1956, vol. 178; Ambrose, E. J., "The Movements of Fibrocytes," Experimental Cell Research, Suppl. 8, 54-73 (1961); Temple, P. A., "Total internal reflection microscopy: a surface inspection technique," Applied Optics, Vol. 20, No. 15, August 1981; and D. Axelrod, in Fluorescence Microscopy of Living Cells in Culture, Part B, ed. D. L. Taylor and Y- L. Wang, (Academic Press, New York, 1989), Chap. 9. Stoney, Ambrose, Temple, and Axelrod disclose optical evanescent light field microscopes in which the light that enters the objective pupil is evanescent field light that has been scattered from a sample surface. However, in all of these microscopes, the sample is transilluminated, with the illumination incident at beyond the critical angle such that the evanescent field from the sample surface is received. It is then necessary that the sample be transparent at optical frequencies, or is made thin enough to be transparent. Scanning devices which rely on scattered evanescent field light are described in: Fischer, U. Ch., Durig, U. T., and Pohl, D. W., "Near-field optical scanning microscopy in reflection," Appl. Phys. Lett., Vol. 52, No. 4, pp. 249-251 25, January 1988. Fischer et al. disclose a near-field optical microscope in which the sample is not transilluminated but is rather illuminated in reflected light. Further, this reflected light is in the form of an evanescent field from a dielectric plate into which light is launched at greater than the critical angle, by means of a coupled prism, so that it undergoes multiple total internal reflections, giving rise to the evanescent field. However, Fisher et al. restrict the evanescent field with an aperture in a metal opaque coating on the total reflection surface of the dielectric plate. This aperture is smaller than the wavelength of light so that an improvement in lateral resolution beyond the normal Abbe limit is achieved, but at the cost of having to scan the aperture relative to the sample to build up an image. A further cost is that energy throughput is very low, making extension to analytical optical techniques such as spectroscopy problematic. Devices which utilize transillumination of transparent samples are described in R. C. Reddick, R. J. Warmack, and T. L. Ferrell, "New form of scanning optical microscopy," Phys. Rev. B 39, 767-70 (1989). Reddick et al. disclose transillumination of thin and transparent samples with evanescent light, but the entrance pupil in Reddick et al is not an objective in the conventional microscopy sense. Rather it is an optical fiber that is scanned over the sample, close to the sample surface. Thus there is a loss of flux throughput, and vertical resolution is limited by the mechanism that controls the vertical position of the fiber relative to the sample. A means of scanning in the xy plane is also required, preventing true real-time whole-field imaging. While the art describes a variety of devices that utilize evanescent field illumination for investigating surface characteristics, there remains a need for improvements that offer advantages and capabilities not found in presently available instruments, and it is a primary object of this invention to provide such improvements. It is another object of the invention to extend the use of photon tunneling imaging to scattering surfaces. Another object is to extend the use of photon tunneling microscopy to rough surfaces by increasing the detectable vertical tunneling range. A still further object is to provide a fixed focus surface profiler requiring little or no skill in microscopy and is simply placed onto the sample and, through kinematic principles and gravity loading, automatically provides a dark field photon tunneling image. Another object is to provide a portable device for on-line web surface roughness inspection and measurement. Another object is to provide a portable device that can be driven about the surface. Another object is to provide a portable surface profiling device that provides position information about the probe relative to the sample surface. Another object is to provide an inverted microscope upon which samples can be easily placed for automatic fixed focus imaging. Another object is to detect light scattered out of the evanescent field by a sample surface in reflected light, rather than in transillumination. Another object is to decouple the illumination optics and optical path to the imaging/collection optics when using the evanescent field to illuminate a surface. Another object is to eliminate the flexible transducer and oil coupling of standard photon tunneling microscopy. Another object is to allow tunneling through water and other materials with refractive index greater than one. Another object is to allow wider fields of view and a larger range of magnifications than standard practice. Another object is to allow greater control over the illumination incident angle, polarization, wavelength, coherence, and phase. Another object is to provide devices for evanescent field illumination that are thin enough to allow use of standard working distance objectives. Another object is to allow photon tunneling microscopy with imaging optics having a numerical aperture equal to or less than unity. A further object is to reduce sensitivity to optical inhomogeneity in the sample surface relative to topography. Another object is to provide real-time 3-D imaging and full statistical analysis. Another object is to provide means for spectral analysis that can be used with photon tunneling for calibration and separation of sample optical inhomogeneity from topography. Other objects of the invention will be obvious, in part, and, in part, will become apparent when reading the detailed description which follows. SUMMARY OF THE INVENTION This invention generally relates to the fields of surface imaging and fine detail measurement through the application of evanescent field illumination and, in particular, to apparatus by which bound evanescent field illumination may be converted by scattering into propagating light that is subsequently collected and imaged for downstream display and metrology purposes. A number of different dark-field, photon tunneling portable probes or optical heads are described along with different ways of providing the necessary evanescent field illumination. In each, the sample that is illuminated by the evanescent field scatters the bound energy in accordance with its scattering properties and the proximity of local surface details to the evanescent field. The scattered light is collected and imaged where the intensity variation of the image varies in accordance with the surface details. Image intensity is correlated with calibrated geometry to provide a means for encoding the intensity variations as variations in gray tones. The intensity of the image is preferably mapped via remote signal processors to a gray tone image in the form of a signal which bears tone density corresponding directly to surface feature dimensions. The gray tone signal thereafter may be used as an input to a number of display possibilities including digital monitors, 3-D oscilloscopes, or other forms of display. In the absence of a sample surface, light is reflected totally outside of the entrance pupil of the imaging system and thus creates a dark field. Hence, the name "dark-field" photon tunneling imaging. In the normal photon tunneling microscope (PTM), the evanescent field is converted into propagating light by refraction into smooth surfaces and away from the entrance pupil of the microscope. In the absence of a sample surface, the light is reflected totally into the microscope entrance pupil, for a bright field. In addition, the dark field PTM is different from normal dark field microscopy in that the illumination is evanescent light, and not propagating light. Sample features get brighter against a dark background as they get closer to the microscope, so height information is obtained. In normal dark field, sample feature brightness depends on their size and slope, rather than their height. In accordance with a further feature of the invention, the imaging optics and illumination optics are decoupled so that advantages are enjoyed in flexibility of light control as well as image magnification and numerical aperture. In accordance with still a further feature of the invention, there is provided by this optical configuration opportunity for novel mechanical and optical configurations and mounting where alignment, proximity placement with the sample surface, and image focus are all automatic, such that the resulting instrument comprises a compact, portable, robust, and easy to use 3-D surface profiler. In accordance with yet another feature of the invention, optical configurations and mounting are provided that facilitate analytical and biological applications. In general, the inventive imaging probes are best used on samples that are scattering and rougher since very smooth surfaces scatter very little light energy into the pupil. On the other hand, these systems are less sensitive to the optical properties of the sample because the conversion of the evanescent field is largely by scattering, rather than by diffraction/refraction. Also, the imaging and illumination optics and optical paths are easily decoupled so that the polarization, phase, incident angle, and wavelength of the illumination is easily controlled. Finally, because the imaging optics are separate and do not require a numerical aperture of one or greater, a much larger range of magnifications, fields of view, and numerical apertures are available commercially. BRIEF DESCRIPTION OF THE DRAWINGS The structure and operation oft he invention, together with other objects and advantages thereof, may best be understood by reading the detailed description in connection with the drawings in which unique reference numerals have been used throughout for each part and wherein: FIG. 1 is a diagrammatic partially elevational and partially perspective view of a dark-field, photon tunneling system of the invention; FIG. 2 is an enlarged, diagrammatic, elevational view of a portion of the system shown in FIG. 1; FIG. 3 is a flow chart illustrating the operation of various embodiments of the invention; FIGS. 4a-4d illustrating how light is scattered or not into the entrance pupil of the detection section of the embodiments of the invention depending on whether or not there is a scattering sample present and the corresponding images that may be observed on the display section; FIG. 5 is a diagrammatic view of an alternate embodiment of the invention in which a microscope looks through an aplanatic sphere to image scattered light; FIG. 6 is a diagrammatic, enlarged elevational view of a portion of the embodiment shown in FIG. 5; FIGS. 7, 8, and 9 are diagrammatic illustrations in which a decoupled light source, objective, and detector camera are arranged for use in another embodiment of a dark field photon tunneling imaging system according to the invention; FIGS. 10 through 24 are diagrammatic views illustrating variations of transmissive body geometries that can be used as the total internal reflection element for the invention; FIGS. 25a through 30 are diagrammatic views of mounting arrangements in which the imaging section a dark field photon tunneling system comprises one pod in a three pod kinematic base for a portable surface profiler; FIG. 31 is a diagrammatic of a viewing section of a dark field photon tunneling system that has been inverted to form yet another embodiment of the invention; FIGS. 32a-32c, 33, and 34 are diagrammatic views showing a dark field photon tunneling microscope mounted as a spring-loaded piston in the center of a kinematic three point ring, in both an inverted as well as upright configuration; FIG. 33 is a diagrammatic illustration of a dark field photon tunneling microscope configured and used as a hand-held device that is placed in contact with the sample; FIG. 34 is a diagrammatic of the dark field photon tunneling technique applied to measuring flying height of a magnetic data recording head in a typical drive arrangement for floppy or rigid disk magnetic media; FIG. 35 is a diagrammatic of a dark field photon tunneling technique applied to high resolution writing and reading of optical data in an optical data drive; FIG. 36 is a diagrammatic view showing how the coupled imaging and illumination optics of a typical photon tunneling microscope can be modified for dark field photon tunneling microscopy; and FIGS. 37 and 38 show representative displays of 3-D and gray scale images, profilometry and roughness statistics such as are available from the dark field photon tunneling microscope DETAILED DESCRIPTION This invention generally relates to the fields of surface imaging and fine detail measurement through the application of evanescent field illumination and, in particular, to methods and apparatus by which bound evanescent field illumination may be converted by scattering into propagating light that is subsequently collected and imaged for downstream display and metrology purposes. A number of different dark-field, photon tunneling imaging probes are described along with different ways of providing the necessary evanescent field illumination in an uncoupled manner. Various other applications relating to magnetic disk flying height measurement, optical surface measurement, probe use, and optical storage are also described. Before proceeding with a description of the inventive probes, dark field photon tunneling will be described with reference to a stationary system. Reference is now made to FIG. 1 which shows a dark field, photon tunneling imaging system of the invention designated generally at 10. As will be seen system 10 is particularly suitable for visualizing and measuring at full field and in real time the microtopographic features of roughened or other surfaces that have small scattering features but otherwise appear smooth. In particular, system 10 and other embodiments of the invention are suitable for use in measuring the surface characteristics of aircraft and space vehicle structures, paint, paper, and fabrics. Because of the lighting and detecting mechanisms employed in dark field imaging, it is possible to decouple the illumination and imaging paths so that probes remote from the imaging sections can be contrived. As seen in FIG. 1, system 10 comprises an illumination and imaging section 12 and an image processing and display section 14. Section 14 may be any well-known general purpose computer or work station 16 having a CPU, RAM memory, hard and floppy drives, input devices such as a keyboard 18 and mouse 20, and color display monitor such as at 22. Preferably, computer 16 has 16 or more megabytes of RAM and is otherwise equipped with high throughput data and video buses. The internal video card is preferably one selected with two or more megabytes of on-board memory and is capable of generating 32-bit or more colors for high tone resolution. Internal signal and image processing programs may be stored on computer 16's internal hard drive and transferred to RAM in the usual way for any processing needs required. As best seen in FIGS. 1 and 2, illumination and imaging section 12 comprises an optical head 24 for illuminating and contacting the surface of a sample to be visualized and measured, and a video camera 26 for collecting and imaging propagating radiation scattered from a sample's surface. Video camera 26 may be any suitable conventionally available type of the desired spatial and tonal resolution. Preferably, video camera 26 is a wide-angle type and may have a magnifying or slightly minifying objective lens 28. Objective lens 28 may alternatively be a zoom lens of appropriate tele to wide angle design. Located at the plane of best focus of objective lens 28 is a photo detector 30, which may be a conventional CCD or vidicon tube. Video signals generated from photo detector 30 may be processed on board camera 26 via a resident chip 32 for that purpose or may be sent via an appropriate board resident in a slot in computer 16. In either event, camera 26 and computer 16 may be configured in well-known manners so that video signals may be digitized to generate digital images that can change for display on monitor 22 at real-time rates or nearly real-time rates. Optical head 24 comprises a bulk optic, prismatic dielectric body 34, having, among others, a light entering facet 42, a sample contacting facet 44, and a light emitting facet 46. Optically coupled to facet 42 is an illumination section comprising a collimating optical section 36 consisting of a tube in which are resident suitable collimating optics in the form of spaced apart lenses 48 and 50. An illumination source 38 is provided, and radiation emitted by source 38 is directed through facet 42 to facet 44. As best seen in FIGS. 2 and 4a, light in transmissive dielectric body 34 is incident at an angle equal to or greater than the critical angle θ c so that it totally reflects at the interface with a less dense medium such as air or water, i.e., at facet 44. It is understood that dielectric body 34 can be transmissive over all or any part of the electromagnetic spectrum, including but not limited to ultra-violet, infrared, or even x-ray and millimeter wave extremes, depending on the application. Although the incident light is shown to be a collimated beam, it need not be, and in fact imaging resolution improves if the incident light is less coherent. As seen in FIG. 4a, an evanescent field 52 arises at the boundary between facet 44 and the lower index medium opposite it (usually air, but this medium may also be water or another low index medium). Evanescent field 52 has an amplitude that decays exponentially with distance from the surface of facet 44. The strength of evanescent field 52 is given by: ##EQU1## where E 0 is the amplitude of the electric field associated with the photon in the medium comprising body 34 and, d p is the penetration depth in the less dense medium at which E 0 decreases to E 0 /e and where: ##EQU2## and λ 1 is the wavelength in the denser medium, θ is the incidence angle, and n 21 is the ratio of denser to lower indices of refraction at the boundary of facet 44. The actual penetration depth, where E 0 falls to the limit of detectability, is dependent on these variables as well as both the photodetector sensitivity and the sample optical properties, and is typically approximately 0.75λ. However, the evanescent field, however small in intensity, can exist sensibly for tens of wavelengths, if the parameters on the last above equation are optimized. As is well-known, evanescent field 52 penetrates normal to the surface of facet 44 to the depth indicated above. Consequently, it extends beyond the physical boundary of facet to a predetermined depth and can be interrupted by a sample placed in close proximity to facet 44. The action of doing so will now be explained with reference to FIGS. 4a to 4d. In FIG. 4a, an incident beam of illumination 54 is shown illuminating facet 44 at the critical angle θ c when no sample is in place. Normal to facet 44 is shown the entrance pupil 56 of video camera 26. Again, evanescent field 52 is shown extending beyond facet 44. Without a sample, the incident beam 54 is totally internally reflected from facet 44 where it is directed as a beam 58 along the indicated path to emerge from facet 60. Since no light enters entrance pupil 56 under these conditions, the image displayed on monitor 22 is completely dark as shown in FIG. 4b at 62. The conditions that prevail when a sample is brought into contact with evanescent field 52 are illustrated in FIG. 4c. Here, a sample 64 with a scattering surface, such as a rough surface, interrupts evanescent field 52 to scatter illumination 68 into entrance pupil 56 where it is used to form an image of the sample surface. The intensity of the scattered radiation, which has been converted from bound energy at the surface of facet 44 to propagating radiation in dielectric body 34, depends on the local scattering properties and the local depth of penetration of the evanescent field of the sample's microtopographic features. The resultant image displayed on monitor 22 is shown in FIG. 4d at 70 where it can be seen that an image area 72 is displayed against an otherwise dark field. Hence the name "dark field" photon tunneling imaging. Because the signal is measured with respect to a dark field, this invention provides generally advantageous signal to noise ratios. In practice, a sample, like 64, is brought into contact with evanescent field 52 with a three-axis, micro manipulable translation stage 40 such as is shown in FIG. 2. Here, a sample 74 is shown sandwiched between facet 44 and the surface of stage 40. Experience has shown that the intensity of propagating radiation converted from the evanescent field can be high enough to visualize and measure surface details ranging from between less than a few angstroms to about 10 micrometers. Useful signal may also be generated with this illumination system to image both reflecting and transmissive materials. FIG. 3 is a general flow chart for the operation of imaging system 10 as well as other embodiments to be described. As shown there, first an evanescent field is formed at a suitable measuring interface as indicated at block 80. Radiation from the bound evanescent field is converted by scattering and proximity into propagating radiation at block 82. The intensity of the propagating radiation depends on the local scattering properties and local depth of penetration of the sample surface into the evanescent field. Following at block 84, the propagating radiant energy is collected and imaged onto a photo detector as by the CCD of video camera 26. An image signal is then generated in block 86 and the intensity of the image signal varies in accordance with the microtopographic features of the sample under scrutiny. The image signal from block 86 is encoded as a gray tone map signal in which gray tones represent feature dimensions. This may be via a look-up table (LUT). The gray tone map signal is then formatted for 2-D and 3-D display via a computer or other suitable video signal processor at block 90. The formatted signal is then displayed at block 92 with monitor 22 or other suitable display as will be subsequently described. The dark field imaging system of the invention is seen to be complementary to the usual photon tunneling microscope (PTM) of the prior art in that scattering surfaces such as paper are difficult to view in the normal PTM because the scatter competes with the gray scale tunneling image, while a smooth surface that is viewed easily with PTM will be largely invisible in dark field PTM because the image in the latter is formed by what is scattered out of the evanescent field, which is very little for such a smooth surface. Another important observation is that in FIG. 4a for the dark field PTM, the entrance pupil 56 is clear of the incident illumination in the dark field PTM, whereas the entrance pupil is the same as the exit pupil for the illumination in the photon tunneling microscope. This separation or decoupling of the illumination and imaging optics in the dark field PTM allows much more freedom in control of the illumination to maximize, for example, the vertical tunneling range that is the sensible amplitude of the exponentially decaying evanescent field, or to control polarization more easily so that non-dielectric samples may be viewed. It should be reiterated at this point that the extent of similarity of dark field PTM to a normal dark field optical microscope ends with the representation of a sample as bright against a dark background, in that the light illuminating a sample in the latter is not evanescent, but is rather the usual propagating electromagnetic field. Therefore, no resolution enhancement is enjoyed in the lateral plane. Nor is there any relation between height and gray scale that only comes with the use of evanescent illumination. Rather, brightness of an object depends on its scattering geometry that includes slope and spatial lateral size, for example. Reference is now made to FIGS. 5 and 6 which illustrate another embodiment for practicing the invention. Here, there is shown generally at 100 a dark field imaging system that comprises an illumination and imaging section 102 and display section 110. Section 102 comprises an optical head 106 which in turn comprises an aplanatic sphere 112 having a slightly curved contact surface 120. Illumination is provided at the critical angle via a fiber optic bundle 116 feed by a source such as a laser, LED, or other solid state device. Suitable radiation controlling filters may also be present in this optical path to control state of polarization or spectral content of illuminating radiation or both. As before, scattered radiation, such as that indicated at 124, propagates to the entrance pupil of an objective lens 122 of a microscope 104. Microscope 104 is focused through aplanatic lens 106 onto surface 120. In this manner, microscope 104 forms an enlarged image of a sample and aplanatic sphere 106 may be used to aide in the correction of spherical, coma, and astigmatism aberrations in the usual way. The image formed by microscope 104 is viewed by a photometric vidicon 108 of conventional design. Vidicon 108 converts the gray scale scattered image into a video signal that is restored by a three-axis oscilloscope 129 to a real-time 3-D image of the sample micro topography while controlling perspective. Other peripherals include a video image processor 122, a VCR 126, a 2-D gray scale image display 124, and an XYZ screen video transfer section 128. An XYZ translation stage 118 is provided for sample positioning as before so that a sample, such as that at 130, can be brought into proximity with the curved surface 120 of aplantic sphere 112. Curved surface 120 facilitates contacting the sample and, as will subsequently be explained, provides other functions. Previously, the inventive dark field photon tunneling imaging system embodiments were intended to be used in stationary modes of operation with the sample moveable with respect to the optical heads. Because the dark field approach permits decoupling between the illumination and imaging optical paths, other optical head configurations are possible, and these will now be taken up. In FIG. 7, an imaging device, a video camera, vidicon, or microscope is equipped with an objective 142. Dark field objective 142, which contains the entrance pupil 56 as shown in FIG. 4a, is made to image through a transmissive dielectric body 144 that corresponds to body 34 in previous figures. A light source 148 is incident on a prism 146 that is coupled optically to dielectric body 144 so that total internal reflection of the light occurs at least at the distal surface 145 of body 144 so that an evanescent field exists outside of distal surface 145. The sample to be imaged is brought proximal to distal surface 145 so that it is within the evanescent field as was illustrated in FIG. 4c. Although the light source 148 is shown as a unidirectional one, and in some cases this type of illumination is preferable, there are many possible illumination configurations as will be illustrated subsequently. In FIG. 8, a more complete dark field PTM instrument 150 is seen to include, in one configuration, an imaging detector such as a CCD camera 152 onto which the real image from an objective 154 is projected. Alternately, a projection eyepiece can be included in this arrangement for increased magnification, and objective 154 can be of even zoom variable magnification. Illumination is provided via a totally internally reflecting body 156 where the evanescent field is formed at a distal end 158. A light source 159 is coupled to body 156 via a bent up section 157. Another useful configuration in shown in FIG. 9. Here, a system 160 uses a coherent imaging fiber optic bundle 162 to relay the image formed by an objective lens 164 to a remote camera and illumination source 166. This embodiment is particularly suitable for use in tight and otherwise inaccessible spaces. A number of additional illumination configurations are possible, as illustrated in FIGS. 10 through 24, and each offers its own unique advantages. The common element in all of these variations is that total reflection of the illumination is made to occur at least at the distal surface of the body in order to cause the necessary evanescent field illumination outside of and adjacent to the distal surface. In FIG. 10, a transmissive body assembly 170 in the form of a figure of rotation provides the optical head functions of previous embodiments. Body assembly 170 includes angled incident entrance facet 172 that extends around its circumference. This allows several, or more than one, discrete light sources (174, 176, and 178) to be placed at any azimuthal angle about the axis of rotation. It should be understood that "light source" or lighting means as used in this disclosure can be a laser diode, a light emitting diode (LED), a fiber bundle carrying light from a remote source, incandescent, halogen, or other light source. At plane 175 can be placed a polarizing filter, wavelength filter, phase plate, or other light controlling device. In FIG. 11, the discrete light sources of FIG. 10 are replaced with an annular light bulb source 180 for lighting the angled circumferential edges 181 of an annular slab block 182 to illuminate the sample from all directions. In FIGS. 12 and 13, light from a source is incident normal to the surface 192 of a flat plate 190 that is parallel to the distal surface 194 but is totally internally reflected by angled edge surface 196 so that it undergoes multiple total internal reflections 198. The light source can be discrete as at 200 or tubular as at 202. In FIG. 12, an optional truncated or "aplanatic" sphere 204 may be made integral with the flat plate 190 in order to improve the image according to aplanatic principles. In FIG. 14, the light source is a monochromater 210 so that analytical dark field PTM can be done. An optional detector 212 is shown at angled edge surface 196 of flat plate 190 which is now inverted. The main detector here would still be an imaging camera that is not shown here but elsewhere. An additional illumination embodiment is illustrated in FIG. 15. As with the embodiments seen previously, the embodiment discussed here would serve as the optical head. In FIG. 15, the transmissive body is in the form of an aplanatic sphere 220 that is truncated at an optically prescribed center thickness 212 such that the image through aplanatic sphere 220 is improved over a parallel flat plate by being aplanatic. The illumination is incident to sphere 220 at angles equal to and greater than the critical angle θ c . The dotted lines 214 indicate an optional metallic coating, the purpose of which is to allow the incident light to enter the sphere at other angles. Surface 216 is the total reflection surface from which the evanescent field emanates. The dotted radius 218 indicates a preferred curved surface to facilitate proximity between the aplanatic sphere and the sample (which is not shown), and also to facilitate sliding of the aplanatic sphere along the sample in the XY plane. The slight curvature does not affect the optical quality of the image to any significant degree, and yet greatly simplifies the practical use of this aplanatic sphere. A complete sphere 230 is shown in FIG. 16. While the imaging quality of a completely spherical transmissive body is not as good as the aplanatic sphere or even the flat plate bodies, it is acceptable for certain applications where a rolling optical probe is needed. The spherical geometry allows for rolling contact with the sample, analogous to the tip of a ball point pen. This embodiment is beneficial where sliding contact between the transmissive body and the sample are undesirable. FIG. 17 shows a rotated solid Dove prism 240 with a curved distal surface indicated by the dotted radius line 244 in comparison to a flat distal surface, with curvature again to facilitate contact with the sample and aid in sliding over the sample in the XY plane. This curvature, the radius of which can be measured and known, also may serve to calibrate height to gray scale in the photon tunneling scattered image. Light enters the Dove prism as a beam 242 as shown. FIG. 18 shows another variation of a transmissive body that is a normal Dove prism 250 except that opposing the incident entrance plane 252 is a truncated plane 254 perpendicular to the distal surface 258, or as indicated by the dotted line 256, parallel to the incident entrance plane 252. This geometry affords advantages in some analytical applications, such as spectroscopic. In FIGS. 19 through 23, the transmissive body is a very thin plate which serves to guide the light that is launched into it at and beyond the critical angle with multiple total internal reflections. It should be understood that in all of these embodiments, the means of launching the light into the light guide can be a single unidirectional means as shown, or a multitude of discrete like means, or a continuous figure of rotation of the means shown for omnidirectional illumination, analogous to the illumination schemes discussed earlier. The advantage common to all of these embodiments is that an objective, such as that shown at 270 in FIG. 23, can now have a very small working distance and higher numerical aperture, thus affording higher lateral resolution and a wider range of magnifications and fields of view, as well as affording a much more compact device overall. As with the other embodiments, it is understood that these thin plate waveguides can substitute for the transmissive slab body assemblies. In FIG. 19, the means of launching light into a waveguide 280 is a prism 282 optically coupled to waveguide 280, so that light incident to prism 282 is at or beyond the critical angle and so undergoes multiple total internal reflections within waveguide 280, giving rise to evanescent field illumination at the distal surface 284. In FIG. 20, the means of launching light into a waveguide 290 is a set of micro prisms 292. This set of micro prisms 292 could also be in the form of periodic microstructure known as diffractive optics. Here, the periodic structures are provided with a spatial period smaller than the wavelength of the light from the lighting means such that the evanescent field is caused by diffraction of light by the periodic structure. In FIG. 21, the means of launching light into a waveguide 300 is a holographic optical element, or HOE 302. In FIG. 22, a waveguide 310 is flexible, and the means of launching light into the flexible waveguide 310 can be a prism 312, HOE, diffractive optics (kinoform), or micro prisms. In FIG. 23a and 23b, a waveguide 320 has lower refractive index core 324 in which the total internal reflections occur. Core 324 is exposed on a distal surface 326 directly to the sample 328, while the cladding 330 above core 324 maintains total reflection in the event of external contacting material. It will be appreciated that a thin transmissive parallel plate or flexible membrane may be used to contact to the distal surface of any of the above transmissive bodies and act as a disposable protective covering. If damaged during sliding contact with the sample, it is removed and replaced, so that expensive damage to the main distal surface is avoided. Hard, diamond or other antiabrasion coating may also be applied in well-known manners to the distal end to avoid damage. Although not shown, all of the transmissive body embodiments discussed above and shown can have a metallized distal surface, with the thickness of the metallization chosen for the illumination wavelength and the incident angle such that a plasmon field is excited on the metal layer by the evanescent field. Therefore, all of the foregoing bodies can be used for surface plasmon microscopy, with all the benefits that brings to, for example, fluorescence microscopy. Similarly, all of the transmissive bodies discussed above can have multiple dielectric thin film coatings on their distal surfaces in order to cause optical resonance, such that the amplitude of the evanescent field is increased for better signal to noise, and better vertical resolution. Also, the entrance and exit faces of the transmissive bodies can be coated and arranged geometrically so that the transmissive body becomes a cavity resonator. And finally, all of the distal surfaces can have non-smooth surfaces, where the random or periodic surface structure has a spatial period smaller than the wavelength of the illuminating light in order to generate a diffraction-born evanescent field for even higher lateral resolution such as with the diffraction grating 331 formed in a pod transmissive body 332 shown in FIG. 24. In FIGS. 25 through 33, several probes including ways of mounting and using the dark field photon tunneling imaging system are illustrated. The many embodiments of the illumination and transmissive bodies discussed above can be used and combined with the mountings to be described. In FIG. 25a, the transmissive body is shown at 340 as one of three pods, the other two being shown at 342 and 344. All three pods are arranged in kinematic fashion on a baseplate 346 such that when baseplate 346 is placed on a reasonably flat sample surface, the distal surface of the transmissive body 340 is automatically in contact with the sample surface, in focus, and aligned. In this way, the dark field PTM functions ergonomically less like a traditional microscope and more as a computer "mouse", except that is a powerful surface profiler. This compact, portable device is easily placed on an airplane wing, or production web of paper, or painted hood of an automobile, for example. In FIG. 25b, the baseplate and pods are shown in profile with the rest of the dark field photon tunneling imaging system components mounted to the baseplate. These include an imaging and illumination section 348 and a control section 350. Imaging and illumination section 348 and control section 350 operate in a manner similar to system 160 of FIG. 9. Many mechanical and other features can be incorporated into the baseplate 346 to provide devices with different capabilities depending on the specific visualization or measurement tasks. Although they will described and shown individually in what follows, the features can in many cases be combined. In FIG. 26, two non-optical pods 360 are motor driven wheels or balls, with one means of drive indicated by motors 362 and 364. The baseplate 346 can then be driven about the sample surface. In FIG. 27, means for air bearing levitation 366 are incorporated into or around all three pods (340, 342, and 344). The air pressure can be used not only for frictionless movement about the sample, but can also control the separation between the distal surface of pod transmissive body 340 and the sample, and with differential control over any two pods, alignment of the optical axis normal to the sample surface is achieved (this is important on non-flat sample surfaces). In FIG. 28, two of the pods (only one shown at 370) are mechanically adjustable in height with micrometers 368 or piezo drives to similarly adjust alignment on non-flat surfaces. In FIG. 29, a counterbalance weight 374 is added such that the two non-optical pods (only one shown at 372) are at the fulcrum point of the device. Sliding weight 374 toward or away from the fulcrum point adjusts the amount of gravity loading of the pod transmissive body 340 against the sample. In FIG. 30, a computer "mouse-like" feature is incorporated into the baseplate 346 so that the position and movement of baseplate 346 and the imaging and illumination sections may be indicated in a well-known manner on a computer screen. For this purpose a ball 376 and well-known positioning technology may be incorporated into baseplate 346. Also in FIG. 30, a marker 378 is added to baseplate 346 so that areas of interest, or defective areas, can be marked in industrial applications. Marker 378 may be any of well-known types such as an inking device. As seen in FIG. 31, the dark field photon tunneling imaging system is very convenient to use also in an inverted position, where a transmissive body distal surface 382 is facing upward. The operator can then hand-place a sample onto the distal surface 382 for immediate images of the sample topography and very rapid throughput of samples. In FIG. 31, an inverted system designated at 380 is accomplished by combination of the portable base 383 similar to the one discussed above nested in a receiver base 385 that has, for example, an optional port for a monochrometer 386 for spectral analysis in addition to topographic profilometry. An optional arm 390 is shown that serves to apply loading to the sample with a device 384 that is an air jet, spring loaded piston, or other mechanism for that purpose. Arm 390 may be raised and lowered into position via a column 392. Other useful, but optional features not shown, are means for heating or cooling the transducer pod and sample, or incorporation of patterning of the illumination for manipulation of microscopic particles. Another useful mounting for a dark field photon maneling imaging system is seen in FIGS. 32a-c. Here, an imager 402 is mounted as a spring loaded piston 414 within a base 400 having three equally spaced pods or balls 404, as is found in an optical spherometer. The distal surface 406 of the transmissive body automatically contacts a flat or curved sample placed on the three pods 404. Piston 414 is mounted in a bored out block 410 in which there is placed a spring 412. In optical fabrication, this mounting configuration provides the optician with both the radius of curvature of the optical surface as well as the roughness statistics and topography, and the degree of optical polish. This configuration can also be used in a downward orientation, as in the center of the base described previously with pods 404 contacting the surface along with the piston 414. In FIG. 33, the dark field photon tunneling imaging and illumination sections are configured as a hand-held probe 420, with the distal surface of the transmissive body 422 placed against the sample 424. The dark field photon tunneling invention can be applied to measure flying height of magnetic read-write heads above floppy or rigid disks in the standard test drives, such as Guzik™, used in the industry. In FIG. 34, light 430 is launched into a glass or polycarbonate analog disk 432 so that it undergoes multiple total internal reflections 433 and therefore induces an evanescent field at both glass to air interfaces. Any of the launch methods previously discussed can be used for this application. Alternately, the edge of disk 432 can be chamfered to form an integral prism face, or can be molded with this prism face to begin with. Even a scattering surface on this edge will launch light into disk 432, though less efficiently and with less control over incident angle. An inside edge 437, nearest to a spindle hub 440, is shown and is the preferred light launch site, but even the outer edge of the disk 432 can serve as the launch site. As a real read/write head (434 or 436) approaches the disk surface 450, it is illuminated by the evanescent field, which is convened into propagating light via scatter by the head, rather than by refraction as in the optical proximity device or photon tunneling microscope. The scattered light enters the dark field PTM objective 438. The typical light-colored ceramic heads will scatter the best, but with a more sensitive detector other less scattering heads may be measured as well. The unadorned glass (or other material) disk will work. However, performance may be improved by blackening, metallization, anti-reflection, or high-low stack dielectric coatings that are selectively applied to strategic areas of the disc or head to improve image contrast and signal to noise, such as one or both of the disk surfaces, or the inside or outside rim of the disk. In addition, a magnetic coating may be added if optically transmissive enough to view through (the near infra-red wavelengths transmit best through such magnetic coatings). Also, a planar waveguide may be created in the disk surface facing the head by infusing chemicals to alter the refractive index. The advantages of using the dark field photon tunneling are all those of photon tunneling but, in addition, the viewing axis is normal to the disk and head surface, thereby eliminating the need for additional optics to view at the critical angle. In addition, it will be understood that it is preferable in the measurement of flying heights with head arrangements for double sided media, it is best to use one real head and one transparent head so the full aerodynamic environment of the read/write system is experienced. Referring now to FIG. 35, it can be seen that dark field photon tunneling, with its high resolution characteristics, can be applied with advantage to the field of optical data storage and retrieval. In FIG. 35, a rotating polycarbonate or other similar optical compact disc 460, or a similar optical card scanned in the XY plane, contains data bits in the form of optical scatter sites 472 on or very near the surface, or nanometer-high topographic bumps 470 on the surface, either of which the evanescent field is sensitive to. A flying total internal reflection head 464 and 468, that can be configured as previously described, serves to illuminate disc 460 in evanescent field light. The information bits convert the evanescent field to scattered light that an objective 466 receives. Alternately, the scatter sites or bumps can be internal to the disk. Into the surface of the disk may be formed a HOE, kinoform, or micro prisms such that, when illuminated normal to the disk surface, only the local area below each element of the micro prism or equivalent HOE is illuminated with evanescent light. As illustrated in FIG. 36, the benefits of dark field photon tunneling imaging or microscopy can be enjoyed in a normal photon tunneling microscope through the use of a conversion which restricts the useful NA of the objective to less than 1.0, and also by reducing the level of specular illumination. Both of these expediencies, while not completely necessary, enhance the performance for dark field purposes. Here, a stop 480 is placed in the epi-illumination path to only allow illumination from numerical aperture 1 to the full extent of the objective, usually 1.25. A complementary annular mask 482 is placed in the imaging optical path to mask out light returning at a numerical aperture greater than one, thus assuring the collection of primarily scattered light from a sample. Reference is now made to FIGS. 37 and 38, which show generally at 500, the viewable area of a monitor or other display device. Within the viewable area, it is possible using well-known techniques to provide different windows to display analytical data, such as statistical data, about the properties of a surface along with image of the surface. For example, box 502 could contain autocorrelation information, 504 could be the image, box 506 and enlargement of the image, and line 508 the surface profile. Similarly, boxes with this type of numerical and graphical or pictorial information can be displayed within boxes as box 512 is within box 510 in FIG. 38. The dark field photon tunneling imaging systems of present invention operate in a novel way because detected light has been scattered out of the evanescent light field by the sample surface. As a result, the advantages of this invention are: increased tunneling range, increased topographic sensing, ability to look at thick and opaque samples, ability to look at rougher surfaces, less sensitivity to optical inhomogeneity in the sample, decoupling of imaging and illumination optics for greater freedom to control illumination and imaging characteristics. Further advantages are real time high speed imaging, and high energy throughput that facilitates optical analytical techniques. It will be understood that the present invention is by no means limited to the particular constructions and methods herein disclosed and/or shown in the drawings, but also comprises any modifications or equivalents within the scope of the claims.

Dark-field imaging probes that employs photon tunneling to visualize and measure submicron features of scattering surfaces at full-field and in real-time or substantially real-time. Each probe is particularly useful for characterizing tiny surface features of reflecting or transmissive materials within the subnanometer to several micron range and is readily adaptable for portable use where it can be employed for characterizing large surfaces such as aircraft wings, fabrics, or papers, or the like. The probes comprise an lighting section for illuminating a surface with an evanescent field and a collecting section positioned with respect to the illuminated surface to channel radiant energy converted from the evanescent field by the process of scattering to unbound energy propagating away from the surface. The intensity of the propagating energy, which varies in accordance with the scattering properties of the surface and the local proximity of surface details to the evanescent field, is imaged preferably at low magnification and wide angle by a following imager which may be a vidicon, digital camera, or other photo detector device. Signals from the imager may be sent to distant system components for subsequent signal processing, display, and metrology purposes.

BACKGROUND OF THE INVENTION [0001] The present information relates to an engine control apparatus, a control method and a control system. [0002] Generally, in a vehicle on which an engine is installed, a starter motor is necessary for starting the engine, and this starter motor is driven by means of a chargeable lead storage battery (hereinafter, the storage battery will be referred to as a battery). In addition, a number of auxiliary machines (electric loads) such as lamps, an air conditioner and electric power window systems which are driven by the battery are installed on the vehicle as electric equipment. Then, in order to recover the battery capacity that has been discharged to start the engine and drive the auxiliary machines, a generator (an alternator) is installed on the vehicle which is driven by the engine to generate electric power so as to charge the battery with electric power so generated. [0003] On the other hand, in the vehicle on which the engine is installed, reducing a fuel consumption relative to a mileage or increasing the fuel economy (hereinafter, referred to as an economical running) constitutes a crucial problem to be solved, and to make this happen, a method is adopted in which the combustion efficiency of the engine is increased, or in order to suppress the wasteful fuel consumption resulting when the vehicle is idled or is under light load, the fuel supply to part of the combustion cylinders is stopped when the vehicle is idled or is under light load. In addition, there occurs a case where the engine is stopped completely for economical running when the vehicle is idled. Furthermore, a vehicle is described in Japanese Patent Publication No. 2000-204995A in which its engine is stopped when the vehicle is idled so as to prevent a battery from becoming flat by controlling an alternator while considering the charged condition of the battery. [0004] Furthermore, a vehicle is described in Japanese Patent Publication NO. 10-153159A in which two batteries of the same type are installed on the vehicle as a battery for starting the engine and a battery for electric equipment and charging power generated by an alternator is properly allocated to the two batteries so as to prevent the batteries from being charged insufficiently to thereby realize the economical running. On the other hand, for an engine having a small number of cylinders, an auxiliary battery such as a lithium ion battery is installed in addition to a normal battery, so that an economical running engine stop is carried out in which the engine stops idling when the vehicle comes to rest. In addition, in an economical running system like this in which two batteries are installed on a vehicle, the applications of a main battery (a lead battery) and an auxiliary battery (a lithium ion battery) are defined as below. [0005] Main battery: to be used as a power supply to the electric loads in such a state that the engine is in operation. [0006] Auxiliary battery: to be used as a power supply to the electric loads in such a state that the engine is not in operation. [0007] Generally, the lead battery having an over 12V output voltage is used for mounting on the vehicle. Meanwhile, the lithium ion batteries having 4V, 8V, 12V, and 16V output voltages respectively are generally used. The lithium ion battery having the 12V output voltage is not suitable as auxiliary battery for mounting on the vehicle. In a case where a lithium ion battery having the 16V output voltage is used as an auxiliary battery, since the output voltage of the lithium ion battery is higher than the output voltage of a main battery, a DC/DC converter becomes necessary which is a voltage regulator for matching the output voltage of the auxiliary battery to the output voltage of the main battery. The DC/DC converter lowers the output voltage of the auxiliary battery when an electric load is driven by the power of the auxiliary battery and raises the output voltage of an alternator when the auxiliary battery is charged with power generated by the alternator. [0008] Normally, a limiter for limiting an output current of the DC/DC converter is incorporated in the DC/DC converter. This limiter is made to be in operation both when the battery is charged while the vehicle is idled and when the battery is charged while the vehicle is driven at constant speed and is made to be released only when the battery is charged while the vehicle is decelerated in order to ensure the output current flowing to the lead buttery when the vehicle accelerate. While a maximum current that the DC/DC converter can supply when the limiter is in operation is, for example, on the order of 15 amperes (hereinafter, referred to as (A), but in the drawings, denoted as A), and a maximum current that the DC/DC converter can supply when the limiter is released is, for example, on the order of 20(A). [0009] In the economical running system on which the two batteries that are configured as has been described above are installed, since a power supply to an electric load during economical running is effected by the lithium ion battery, in the event that the voltage of the lithium ion battery is low, the economical running is prohibited until the lithium ion battery is charged to a high voltage. The charging of the lithium ion battery is carried out by the alternator and the DC/DC converter. [0010] In the twin-battery type economical running system that is configured as has been described above, however, when the lithium ion battery is charged, since the limiter of the DC/DC converter is released only when the battery is charged while the vehicle is decelerated, current that charges the lithium ion battery is not much, leading to a problem that it takes much time to charge the lithium ion battery. Due to this, when the vehicle comes to be idled in such a state that the voltage of the lithium ion battery is low, a long time is needed until an economical running is permitted which is enabled when the voltage of the lithium ion battery is high, leading to a problem that the number of times of permitting the economical running is reduced. SUMMARY OF THE INVENTION [0011] It is therefore an object of the invention to provide an engine control apparatus, a control method and a control system for the twin-battery type economical running system which can restore the voltage of a second battery to a voltage at which an economical running can be permitted as soon as possible while preventing the voltage of a first battery as the main battery from voltage reduction in the event that the voltage of the second battery is reduced by reviewing the charging method of the second battery and the control method of a current limiter of a voltage regulator. [0012] In order to achieve the above-mentioned object, according to the invention, there is provided an engine control apparatus adapted to be mounted on a vehicle including: a first buttery, supplying power to an electrical load of the vehicle when an engine is in a normal operation; a second buttery, supplying power to the electrical load of the vehicle when the engine is in an economical operation; a voltage regulator regulating voltages of the first buttery and the second buttery; and a current limiter, incorporated in the voltage regulator and limiting a current value flowing from the voltage regulator to the second battery to a limiting value, the engine control apparatus comprising: [0013] an engine stopping unit, stopping the engine when a first predetermined condition is established; [0014] an activator, starting the engine when a second predetermined condition is established; [0015] a current detecting unit, detecting the current value when the second battery is charged; and [0016] a changing unit, changing the limiting value between a first limiting value and a second limiting value which is higher than the first limiting value according to the detected current value. [0017] According to the invention, there is also provided an engine control apparatus adapted to be mounted on a vehicle including: a first buttery, supplying power to an electrical load of the vehicle when an engine is in a normal operation; a second buttery, supplying power to the electrical load of the vehicle when the engine is in an economical operation; and a voltage regulator, regulating voltages of the first buttery and the second buttery; the engine control apparatus comprising: [0018] a voltage detecting unit, detecting the voltages of the first buttery and the second buttery; and [0019] switching unit, switching a target voltage value of the second buttery from a first voltage value to a second voltage value which is higher than the first voltage value when the second battery is charged in a case where the voltage of the second battery is no more than a second predetermined value and the voltage of the first battery is no less than a first predetermined value. [0020] The engine control apparatus may further comprises: [0021] a voltage detecting unit, detecting the voltages of the first buttery and the second buttery; and [0022] switching unit, switching a target voltage value of the second buttery from a first voltage value to a second voltage value which is higher than the first voltage value in a case where the voltage of the second battery is no more than a second predetermined value and the voltage of the first battery is no less than a first predetermined value. [0023] According to the invention, there is also provided an engine control method for an vehicle adapted to include; a first buttery, supplying power to an electrical load of the vehicle when an engine is in a normal operation; a second buttery, supplying power to the electrical load of the vehicle when the engine is in an economical operation; a voltage regulator regulating voltages of the first buttery and the second buttery; and a current limiter, incorporated in the voltage regulator and limiting a current value flowing from the voltage regulator to the second battery to a limiting value, the engine control method comprising: [0024] stopping the engine when a first predetermined condition is established; [0025] starting the engine when a second predetermined condition is established; [0026] detecting the current value when the second battery is charged; and [0027] changing the limiting value between a first limiting value and a second limiting value which is higher than the first limiting value according to the detected current value. [0028] According to the invention, there is also provided an engine control system adapted to be mounted on a vehicle including: a first buttery, supplying power to an electrical load of the vehicle when an engine is in a normal operation; a second buttery, supplying power to the electrical load of the vehicle when the engine is in an economical operation; a voltage regulator regulating voltages of the first buttery and the second buttery; and a current limiter, incorporated in the voltage regulator and limiting a current value flowing from the voltage regulator to the second battery to a limiting value, the engine control apparatus comprising: [0029] an engine stopping unit, stopping the engine when a first predetermined condition is established; [0030] an activator, starting the engine when a second predetermined condition is established; [0031] a current detecting unit, detecting the current value when the second battery is charged; and [0032] a changing unit, changing the limiting value between a first limiting value and a second limiting value which is higher than the first limiting value according to the detected current value. [0033] In the related art economical running system employing two batteries, if the voltage of the second battery is lower than the second predetermined value, the function of charging the second battery is not performed except the case of decreasing the speed of a vehicle so that the engine for economical running is not stopped. Therefore, in the related art economical running system employing two batteries, the gas mileage may not be improved. However, according to the present invention, even if a voltage of the second battery is lower than the second predetermined value, the function of charging the second battery is performed under the conditions other than the case of decreasing the speed of a vehicle. Therefore, the voltage of the second battery may be restored in an early stage so that the engine for economical running may be stopped and accordingly, the gas mileage may be improved. BRIEF DESCRIPTION OF THE DRAWINGS [0034] The above objects and advantages of the present invention will become more apparent by describing in detail preferred exemplary embodiments thereof with reference to the accompanying drawings, wherein: [0035] FIG. 1 is a configuration diagram which shows an embodiment of a configuration including an engine and an auxiliary machine as electric equipment to which the invention is applied; [0036] FIG. 2 a is an explanatory diagram which shows an charge control permissive voltage and an economical running prohibitive voltage of a lead battery, and FIG. 2 b is an explanatory diagram which explains the meaning of a voltage in a lithium ion battery; [0037] FIG. 3 a is a flowchart of a first embodiment of the invention which shows an embodiment of a procedure for a current limit control by a DC/DC converter, and FIG. 3 b is a time chart which shows a transition of a current limit value of the DC/DC converter according to the control procedure in FIG. 3 a; [0038] FIG. 4 a is a flowchart of a second embodiment of the invention which shows an embodiment of a procedure for an output voltage control of a DC/DC converter, and FIG. 4 b is a time chart which shows a transition of a output voltage of the DC/DC converter according to the control procedure in FIG. 4 a; [0039] FIG. 5 is a flowchart of a third embodiment of the invention which shows an embodiment of a procedure for controlling the charge of the lithium ion battery; [0040] FIG. 6 is a time chart of a fourth embodiment of the invention which shows transitions of charge current to the lithium ion battery and a current limit value; [0041] FIG. 7 is a flowchart which shows a control procedure shown in the time chart in FIG. 6 ; [0042] FIG. 8 a is a flowchart which shows a first mode of a control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged in a fifth embodiment of the invention and FIG. 8 b is a time chart which shows transitions of the voltage value of the lithium ion battery and the output voltage of the DC/DC converter in the control shown in FIG. 8 a; [0043] FIG. 9 a is a flowchart which shows a second mode of a control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged in a sixth embodiment of the invention and FIG. 9 b is a time chart which shows transitions of the voltage value of the lithium ion battery and the output voltage of the DC/DC converter in the control shown in FIG. 9 a , and FIG. 9 c is a time chart which shows a transition of the output voltage of the DC/DC converter when the voltage value of the lead battery is low in the control shown in FIG. 9 a; [0044] FIG. 10 a is a flowchart which shows a third mode of a control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged in a seventh embodiment of the invention and FIG. 10 b is a time chart which shows transitions of the voltage value of the lithium ion battery and the output voltage of the DC/DC converter in the control shown in FIG. 10 a; [0045] FIG. 11 is a time chart of an eighth embodiment of the invention which shows a fourth mode of a control procedure of the output voltage of the DC/DC converter when the lithium ion battery is charged; [0046] FIG. 12 is a flowchart of a ninth embodiment of the invention which shows a procedure in which the control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged is combined with a limiter control of the DC/DC converter; [0047] FIG. 13 is a flowchart of a tenth embodiment of the invention which shows an embodiment of an output voltage control of the DC/DC converter when a lead battery is discharged; [0048] FIG. 14 is a time chart of the tenth embodiment of the invention which shows transitions of an integrated value of discharge current from the lead battery, the output voltage of the DC/DC converter and the voltage value of the lithium ion battery in the control procedure shown in FIG. 13 ; [0049] FIG. 15 is a time chart of an eleventh embodiment of the invention which explains a correction process in a charge control of the lithium ion battery; [0050] FIG. 16 is a flowchart diagram illustrating an operation of controlling the cancellation of a limiter value in the process of controlling the charge of the lithium ion battery according to a twelfth example embodiment of the present invention; and [0051] FIG. 17 is a flowchart diagram illustrating an operation of converting an output voltage of a DC/DC converter in the process of controlling the charge of the lithium ion battery according to a thirteenth example embodiment of the present invention. DETAILED DESCRIPTION OF THE EMBODIMENTS [0052] Hereinafter, a mode for carrying out the invention will be described in detail based on specific embodiments while referring to the accompanying drawings. [0053] FIG. 1 is such as to show the configuration of an embodiment of an automatic engine stop/start control apparatus of the invention, and an engine 6 is such as to be installed on a vehicle. In addition, electric equipment is an auxiliary machine 4 that is installed on the vehicle, and an electric load such as a starter motor 7 for starting the engine 6 , lamps or an air conditioner, or a motor for opening and closing a window glass which uses a lead battery 1 as its power supply corresponds to the auxiliary machine 4 . The lead battery 1 is connected to the auxiliary machine 4 when an ignition switch SW 1 is on (in an IG position) and is also connected to an alternator 5 which generates power with which the lead battery 1 is charged. An ammeter 1 A, a voltmeter 1 V and the ignition switch SW 1 are provided on a circuit between the lead battery 1 and the auxiliary machine 4 . In addition, a switch SW 2 is provided between the lead battery 1 and the alternator 5 , and the starter 7 is made to start the engine 6 automatically when the ignition switch SW 1 is put in a starter position and a starter switch SW 3 is on. An economical running ECU can start the engine 6 automatically by causing the starter switch SW 3 to be switched on. [0054] On the other hand, in addition to the lead battery 1 , a lithium ion battery 2 is installed on this vehicle. In general, since the output voltage of the lead battery 1 is over 12V, while the output voltage of the lithium ion battery 2 is 16V, the lithium ion battery 2 cannot be connected to the lead battery 1 in parallel as they are. Then, a DC/DC converter 3 is connected to the lithium ion battery 2 , and this DC/DC converter 3 lowers the output voltage of 16V to 13V which is slightly higher than the output voltage of the lead battery 1 when power is supplied from the lithium ion battery 2 and raises the output voltage of 12V of the alternator 5 to 16V when the lithium ion battery 2 is charged by the alternator 5 . In this embodiment, an ammeter 2 A and a voltmeter 2 V are provided at an output of the lithium ion battery 2 , and a switch SW 5 is provided to an earth side of the lithium ion battery 2 . In addition, the switch SW 5 provided to the earth side of the lithium battery 2 is made to be switched on and off by the economical running ECU 10 . [0055] On the other hand, in the twin-power-supply installed economical running system on which the lead battery 1 and the lithium ion battery 2 are installed as has been described above, the starter 7 is started by the lead battery 1 for engine start (with an ignition key) except for cold start and start from an engine stop for economical running, and when the engine is started from the economical running engine stop state, the starter 7 is started by the lithium ion battery 2 . Due to this, in this embodiment, the starter motor 7 is connected to the lead battery 1 and the lithium ion battery 2 via a starter power supply changeover switch SW 4 . [0056] In addition, in this embodiment, an ECU for electronically controlling fuel injection (in the figure, described as EFI-ECU) 9 , the economical running ECU 10 for stopping and starting the engine 6 for economical running and a battery ECU 11 correspond to control units for driving the engine 6 and the auxiliary machine 4 which is electric equipment. The engine 6 installed on the vehicle is controlled by the EFI-ECU 9 and is supplied with fuel from a fuel supply system 8 . In addition, the economical running ECU 10 is connected to the EFI-ECU 9 via a bus 16 , is connected to the DC/DC converter 3 via a bus 12 and is connected to the battery ECU 11 via a bus 13 . [0057] A current limiter 17 is incorporated in the DC/DC converter 3 which can limit the value of current which flows through the DC/DC converter 3 by an external signal. In addition, in this embodiment, five input terminals DDIG, DDON, DDDR, DDRE, and BT are provided on the DC/DC converter 3 to which signals from the economical running ECU 10 are inputted. [0058] The input terminal DDIG is a power supply terminal, to which power supply for operating the DC/DC converter 3 is inputted. A signal which determines whether or not the DC/DC converter 3 is activated is inputted to the input and output terminal DDON. When an ON signal is inputted to the input and output terminal DDON, the DC/DC converter 3 starts to operate, whereas an OFF signal is inputted thereto, the DC/DC converter 3 is rendered inoperable. A signal which determines the output direction of the DC/DC converter 3 is inputted to the input terminal DDDR, and depending on a signal so inputted, it is determined whether current is outputted in a direction where the lead battery 1 exists or in a direction where the lithium ion battery 2 exists. A signal which determines on the operation of the incorporated current limiter 17 is inputted to the input terminal DDRE. For example, when an ON signal is inputted to the input terminal DDRE, the current limiter 17 is allowed to output a maximum current of 50(A), whereas when an OFF signal is inputted thereto, current flowing through the current limiter 17 is restricted, and hence, the current limiter 17 is only allowed to output a maximum current of 15(A). A signal which determines on the magnitude of voltage that is outputted from the DC/DC converter 3 is inputted to the input terminal BT. [0059] The battery ECU 11 detects a state of the battery from detection values of the ammeters 1 A, 2 A and the voltmeters 1 V, 2 V so as to perform charging and discharging controls over the lithium ion battery 2 and the lead battery 1 in cooperation with the EFI-ECU 9 and the economical running ECU 10 . For example, in a case where the charge rate of the lithium ion battery 2 is less than 50%, the battery ECU 11 outputs a signal to terminal DDDR of the DC/DC converter 3 through the economical running ECU 10 so that the current output from the DC/DC converter 3 is oriented to the lithium ion battery 2 . In addition, the economical running ECU 10 controls the switches SW 2 to SW 5 except for the ignition switch SW 1 to be on and off. [0060] In this embodiment, while the engine is stopped for economical running, power is not supplied to the auxiliary machine 4 from the lead battery 1 but is supplied thereto from the lithium ion battery 2 in order to prevent the deterioration of the lead battery 1 . As this occurs, the switch SW 5 is switched on by the economical running ECU 10 , and the output voltage of the lithium ion battery 2 is applied to the ignition switch SW 1 via the DC/DC converter 3 . In addition, since a voltage from the lithium ion battery 2 via the DC/DC converter 3 is higher than the output voltage of the lead battery 1 , power is not supplied from the lead battery 1 to the auxiliary machine 4 in this state but is supplied thereto from the lithium ion battery 2 . [0061] In addition, when the economical running engine stop is ceased to restart the engine, the starter power supply changeover switch SW 4 is switched to the lithium ion battery side, so that power is supplied from the lithium ion battery 2 directly to the starter motor 7 not via the DC/DC converter 3 but through the starter power supply changeover switch SW 4 . On the other hand, power that the DC/DC converter 3 can supply is limited, and when the auxiliary machine 4 is driven by the lithium ion battery 2 , in the event that the auxiliary machine 4 is put under a highly loaded state and comes to require more current than the maximum current supply capacity (50 amperes) of the DC/DC converter 3 , since a required power supply to the auxiliary machine 4 cannot be available only from the lithium ion battery 2 , the lead battery 1 is also used in parallel with the lithium ion battery 2 to meet the required power supply. [0062] FIG. 2 ( a ) illustrates a charge control permissive voltage V 1 of the lead battery, at which the lithium ion battery comes to be chargeable, and an economical running prohibitive voltage VX of the lead battery. FIG. 2 ( b ) illustrates a voltage in terms of controlling the lithium ion battery. The operation of controlling the charge of the lithium ion battery may be performed when the voltage of the lead battery is greater than or equal to a voltage V 1 . In addition, when the voltage of the lead battery is less than or equal to the economical running prohibitive voltage VX, the operation of economical running is prohibited. For example, the value of the charge control permissive voltage V 1 of the lead battery, at which the lithium ion battery comes to be chargeable, is approximately 12 volts. The value of the economical running prohibitive voltage VX is approximately 11.8 volts in general. However, the value of the economical running prohibitive voltage VX may vary according to temperature. [0063] Further, voltage values V 3 , V 4 , V 0 and VZ in order of the magnitude are assigned to the lithium ion battery. The voltage value V 3 is the maximum value of a target voltage to which the lithium ion battery is to be charged. For example, the voltage value V 3 may correspond to 16.0 volts. The voltage value V 4 is an ordinary voltage value of the target voltage to which the lithium ion battery is to be charged. The voltage value V 0 is for performing a control operation for converting a limiter value (in other words, a voltage value for performing a control operation for charging the lithium ion battery). For example, the voltage value V 0 may correspond to 14.8 volts. In addition, the voltage value VZ is the economical running prohibitive voltage VX. For example, the economical running prohibitive voltage VZ may correspond to 12.5 volts. [0064] In the following description, the voltage of the lead battery 1 is denoted by PbV, current which flows from or which is charged to the lead battery 1 is by PbA, the voltage of the lithium ion battery 2 is by LiV, and current which flows from or which is charged to the lithium ion battery 2 is by LiA. In addition, in the drawings, there are cases where the lead battery 1 is briefly described as Pb and the lithium ion battery as Li. Furthermore, in the drawings, ampere (A) which is the unit of current is simply described as A. [0065] In the configuration shown in FIG. 1 , the economical running which is effected by stopping the engine 6 by means of the EFI-ECU 9 , the economical running ECU 10 and the battery ECU 11 will be executed in the following procedure. When the economical running is effected, the economical running is executed when economical running conditions are met completely such as the vehicle is idled while stopped. [0066] When the economical running is performed, the voltage PbV and the current PbA of the lead battery 1 , and the voltage LiV and the current LiA of the lithium ion battery 2 are detected. When the voltage LiV of the lithium ion battery 2 is greater than or equal to the economical running prohibitive voltage VZ shown in FIG. 2 ( b ) and the voltage PbV of the lead battery 1 is greater than or equal to the economical running prohibitive voltage VX, economical running is performed so that an engine 6 may be stopped. When the voltage LiV of the lithium ion battery 2 is less than the economical running prohibitive voltage VZ, or the voltage PbV of the lead battery 1 is less than the economical running prohibitive voltage VX, economical running is not performed. [0067] Since power is supplied to the electric load by the lithium ion battery 2 while the economical running is executed as has been described above, in the event that the economical running is repeated several times, the voltage of the lithium ion battery 2 is dropped to be lower than the charge control execution voltage V 0 , approaching the economical running prohibitive voltage VZ. As this occurs, the lithium ion battery 2 needs to be charged with good efficiency. The invention is such that the charge control of the lithium ion battery is carried out in a case where the voltage LiV of the lithium ion battery 2 is less than or equal to the voltage of V 0 , and the voltage PbV of the lead battery 1 is greater than or equal to the charge control permissive voltage V 1 , so as to control the lithium ion battery 2 to be charged with good efficiency, and a mode for carrying out the invention will be described below based on several embodiments. Control procedures which will be described using flowcharts are to be executed every predetermined period of time. [0068] FIG. 3 ( a ) is such as to show a first embodiment of the invention and shows specifically a current limitation control by the DC/DC converter in the charge control on the lithium ion battery. In this embodiment, the current limit value of the current limiter incorporated in the DC/DC converter is changed according to a value of current which flows into the side of the lithium ion battery 2 while the lithium ion battery is charged, with the voltage value LiV of the lithium ion battery being equal to or smaller than V 0 and the voltage value PbV of the lead battery 1 being equal to or larger than V 1 . In addition, FIG. 3 ( b ) is a time chart which shows a transition of the current limit value of the DC/DC converter according to a control procedure shown in FIG. 3 ( a ). [0069] Note that in the embodiment shown in FIG. 3 ( a ), an example will be described in which the current limit value of the current limiter incorporated in the DC/DC converter is released in the event that the value of current which flows into the side of the lithium ion battery exceeds a reference current value. Namely, an example will be described in which a normal current limit value of 15(A) of the DC/DC converter is released so as to allow current to flow up to a maximum current value of 50(A) of the DC/DC converter. [0070] In step 301 , whether or not the engine is being idled is determined, and if it is determined that the engine is being idled, proceed to step 303 , whereas if it is determined that the engine is not being idled, proceed to step 302 , where whether or not the vehicle is being operated in a constant speed running state is determined. Then, if the vehicle is not in the constant speed running state, this routine ends. If the vehicle is in the constant speed running state, however, proceed to step 303 . In step 303 , whether or not the lithium battery is being charged is determined, and if the battery is not being charged, this routine ends. If the lithium battery is being charged, however, proceed to step 304 . [0071] In step 304 , a value of charge current LiA to the lithium ion battery is read from the sensor 2 A shown in FIG. 1 , and in the following step 305 , whether or not the value of charge current LiA to the lithium ion battery so read is equal to or larger than a reference current value I 0 is determined. Then, if LiA≧I 0 , proceed to step 306 , where the normal current limit value of 15(A) of the DC/DC converter is released, and a charge control on the lithium ion battery is carried out. Namely, the current limit value of the limiter incorporated in the DC/DC converter, which is described in FIG. 1 , is raised from the normal current limit value of 15(A) to the current limit value of 50(A) of the invention to execute the control. This state is shown in FIG. 3 ( b ). [0072] On the contrary, if the determination in step 305 results in LiA<I 0 , proceed to step 307 , the charge control on the lithium ion battery is carried out while the normal current limit value of the DC/DC converter is left as it is. When the steps 306 , 307 are completed, this routine ends. In the first embodiment, since an output voltage of the DC/DC converter resulting when the predetermined condition is established while the lithium ion battery is charged becomes larger than one resulting while the lithium ion battery is charged normally, the lithium ion battery can be charged completely in a short period of time compared to the related-art case. As a result, the frequency of occurrence of the economical running is increased to thereby improve the fuel economy. [0073] FIG. 4 ( a ) is such as to show a second embodiment of the invention and shows specifically an output voltage control by the DC/DC converter in the charge control on the lithium ion battery. In this embodiment, a target charge voltage value TLiV of the lithium ion battery is changed from the normal target voltage V 4 to the target voltage maximum value V 3 or the output voltage maximum value V 3 of the DC/DC converter so as to reduce the charging time of the lithium ion battery while the lithium ion battery is charged, with the voltage value LiV of the lithium ion battery 2 being equal to or larger than the charge control execution voltage V 0 and with the voltage value PbV of the lead battery 1 being equal to or larger than the charge control permissive voltage V 1 . In Addition, FIG. 4 ( b ) is a time chart which shows a transition of the output voltage of the DC/DC converter according to a control procedure shown in FIG. 4 ( a ). [0074] In step 401 , the voltage value LiV of the lithium ion battery and the voltage value PbV of the lead battery are read. In the following step 402 , whether or not the voltage value LiV of the lithium ion battery 2 is equal to or smaller than V 0 is determined, and if LiV≦V 0 , proceed to step 403 , whereas if LiV>V 0 , proceed to step 407 . In step 403 , whether or not the voltage value PbV of the lead battery 1 is equal to or larger than V 1 is determined. If PbV≧V 1 , proceed to step 404 , and if PbV<V 1 , proceed to step 407 . [0075] In step 404 , whether or not the engine is being operated in an idling state is determined, and if the engine is in the idling state, proceed to step 406 , whereas if the engine is not in the idling state, proceed to step 405 , where whether or not the vehicle is in the constant speed running state is determined. Then, if the vehicle is not in the steady speed running state, proceed to step 407 , whereas if the vehicle is in the steady speed running state, proceed to step 406 . In step 406 , the output voltage of the DC/DC converter is changed from the normal output voltage value V 4 to the output voltage value of the invention, that is, the output voltage maximum value V 3 of the DC/DC converter so as to carry out the charge control on the lithium ion battery. This state is shown in FIG. 4 ( b ). [0076] On the other hand, in step 407 to which the routine proceeds if LiV>V 0 in step 402 , if PbV<V 1 in step 403 or if the vehicle is determined not to be in the steady speed running state in step 405 , the charge control is carried out on the lithium ion battery while the output voltage of the DC/DC converter remains at the normal output voltage value V 4 . When the steps 406 , 407 are completed, this routine ends. In the second embodiment, since, an output voltage of the DC/DC converter resulting when the predetermined condition is established while the lithium ion battery is charged becomes larger than one resulting while the lithium ion battery is charged normally, the target charge voltage value TLiV of the lithium ion battery becomes large, and hence, the lithium ion battery can be charged completely in a short period of time compared to the related-art case. As a result, the frequency of occurrence of the economical running is increased to thereby improve the fuel economy. [0077] FIG. 5 is such as to show a third embodiment of the invention and shows specifically a control resulting from a combination of the first and second embodiments in the charge control on the lithium ion battery. In this embodiment, firstly, the output voltage control by the DC/DC converter of the second embodiment which is described by reference to FIG. 4 ( a ) is carried out, and thereafter, the current limitation control by the DC/DC converter of the first embodiment which is described by reference to FIG. 3 ( a ) is carried out. In the third embodiment, since a control procedure from step 401 to 407 is completely identical with the control procedure described in FIG. 4 ( a ), like step numbers will be imparted to like steps, so as to omit the description thereof. [0078] While in the second embodiment, the routine ends when the steps 406 , 407 are completed, in the third embodiment, when steps 406 , 407 are completed, the routine proceeds to step 408 . In step 408 , a value of charge current LiA to the lithium ion battery is read from the sensor 2 A shown in FIG. 1 , and in the following step 409 , whether or not the value of charge current LiA to the lithium ion battery so read is equal to or larger than the reference current value I 0 is determined. [0079] Then, if LiA≧I 0 , proceed to step 410 , where the normal current limit value of 15(A) of the DC/DC converter is released, and a charge control on the lithium ion battery is carried out. Namely, the current limit value of the limiter incorporated in the DC/DC converter, which is described in FIG. 1 , is raised from the normal current limit value of 15(A) to the current limit value of 50(A) of the invention to execute the control. On the contrary, the determination in step 409 results in LiA<I 0 , proceed to step 411 , where the charge control on the lithium ion battery is carried out while the normal current limit value of 15(A) of the DC/DC converter is left as it is. When the steps 410 , 411 are completed, this routine ends. [0080] In the third embodiment, since, when the predetermined condition is established while the lithium ion battery is charged, the value of a voltage outputted then from the DC/DC converter becomes large and the value of a current which flows into the lithium ion battery through the DC/DC converter becomes large compared to those resulting while the lithium ion battery is charged normally, the lithium ion battery can be charged completely in a short period of time compared to the related-art case. As a result, the frequency of occurrence of the economical running is increased to thereby improve the fuel economy. [0081] While, in the first embodiment, the example is described in which the current limit value of the current limiter of the DC/DC converter is released when the value of current which flows through the DC/DC converter becomes equal to or larger than I 0 , in reality, in the event that a state where the value of current which flows through the DC/DC converter equals the normal current limit value I 0 continues for a predetermined period of time or longer, the current limit value of the current limiter of the DC/DC converter is released, and in the event that a state where the value of current which flows through the DC/DC converter lowers below the normal current limit value I 0 continues for a predetermined period of time or longer, the current limiter of the DC/DC converter is restored to operate its normal limiter control. This control will be made to constitute a fourth embodiment of the invention and will be described using time charts in FIGS. 6 ( a ), 6 ( b ) and a flowchart in FIG. 7 . [0082] FIG. 6 ( a ) is such as to show a sensor value which indicates an actual current detected by the current sensor as flowing through the DC/DC converter, the normal limitation value of the current limiter of the DC/DC converter and the maximum current value of the DC/DC converter which results when the limiter is released. In addition, FIG. 6 ( b ) is such as to show a transition of the voltage of the lithium ion battery while the lithium ion battery is charged. In this embodiment, a case will be described where the lithium ion battery is started to be charged at a time ta. [0083] When the charging of the lithium ion battery is started at the time ta, the value of current flowing through the DC/DC converter gradually increases and reaches, at a time tb, the current value I 0 which is the normal current limit value of the current limiter. The value of this current value I 0 is 15(A), for example. In the event that this state continues by a duration of T 0 (ms), the current limit value of the limiter of the DC/DC converter is released at a time tc which is a time resulting after the duration of T 0 has elapsed since the time tb. As a result, current up to a maximum current I 1 can flow through the DC/DC converter. The value of the maximum current I 1 is 50(A), for example. [0084] When the current limit value of the current limiter of the DC/DC converter is released at the time tc, the value of current flowing through the DC/DC converter increases. Thereafter, while the value of current flowing through the DC/DC converter continues to exceed the current limit value I 0 for a certain period of time, thereafter, the current value lowers below the current limit value I 0 at a time td. When this state continues by a duration of time T 1 (ms), the current limiter of the DC/DC converter is restored to operate its normal limiter control at a time te which is a time resulting when the duration of time T 1 has elapsed since the time td, and the current limit value I 0 is set. [0085] Assuming that thereafter, the value of current flowing through the DC/DC converter increases again and reaches, at a time tf, the current value I 0 which is the normal current limit value of the current limiter and this state continues by the duration of T 0 (ms). Then, the current limit value of the limiter of the DC/DC converter is released at a time tg which is a time resulting when the duration of time T 0 has elapsed since the time tf, whereby the state is restored where current up to the maximum current I 1 can flow through the DC/DC converter again. [0086] When the current limit value of the current limiter of the DC/DC converter is released at the time tg, the value of current flowing through the DC/DC converter increases. Assuming that thereafter, while the value of current flowing through the DC/DC converter continues to exceed the current limit value I 0 for a certain period of time, thereafter, the value of current flowing through the DC/DC converter lowers slightly below the current limit value I 0 at a time th. When this state continues by the duration of time T 1 (ms) without any change, the current limiter of the DC/DC converter is restored to operate its normal limiter control at a time ti which is a time resulting when the time T 1 has elapsed since the time th. [0087] Incidentally, when the voltage of the lithium ion battery which has continued to be charged sine the time ta reaches a control termination voltage V 5 at a time tH which is after the time th but before the time ti, the current limiter of the DC/DC converter is restored to operate its normal limiter control at the time tH, and the current limit value I 0 is set. Thereafter, the value of current flowing through the DC/DC converter lowers below the current limit value I 0 , and the voltage of the lithium ion battery comes to reach the normal target voltage V 4 at a time t 1 . [0088] FIG. 7 shows a flowchart which shows a control procedure of the current limiter of the DC/DC converter in a case where current flowing through the DC/DC converter varies as shown in FIG. 6 while the lithium ion battery is charged. Also in this embodiment, the charging of the lithium ion battery is carried out while the vehicle is in the idling state or in the steady speed running state. [0089] In step 701 , the value of current LiA which flows into the lithium ion battery, an actual voltage value LiV of the lithium ion battery and the target charge voltage TLiV of the lithium ion battery are read. In the following step 702 , whether or not the current limit value of the current limiter of the DC/DC converter is 15(A) is determined, and if the current limit value is 15(A), proceed to step 703 , whereas if not, proceed to step 707 . [0090] In step 703 , whether or not the value of current LiA which flows into the lithium ion battery is smaller than 15(A) is determined, and if LiA<15(A), this routine ends. If LiA≧15(A) (in reality, since the current limit value of the current limiter is 15(A), whether or not LiA=15(A) may be determined), however, proceed to step 704 . In step 704 , time T during which LiA=15(A) is counted. [0091] In the following step 705 , whether or not the time T counted in step 704 becomes equal to or longer than a predetermined period of time is determined, and if T<T 0 , this routine ends. If T≧T 0 , however, proceed to step 706 . In step 706 , the current limit value of 15(A) of the current limiter is released, and the counted value of the time T is cleared, the routine ending. As a result, the current limit value of the current limiter is released, current up to 50(A) comes to flow through the DC/DC converter. The procedure from step 703 to step 706 corresponds to the operation occurring from the time tb to the time tc or from the time tf to the time tg in the time chart shown in FIG. 6 . [0092] On the other hand, in step 707 to which the routine proceeds if the current limit value is determined not to be 15(A) in step 702 , whether or not the value of current LiA which flows into the lithium ion battery is smaller than 15(A) is determined, and if LiA≧15(A), proceed to step 711 . If LiA<15(A), however, proceed to step 708 . In step 708 , time T during which LiA<15(A) is counted. [0093] In the following step 709 , whether or not the time T counted in step 708 becomes equal to or longer than a predetermined period of time T 1 is determined, and if T<T 1 , proceed to step 711 , whereas if T≧T 1 , proceed to step 710 . In step 710 , the current limit value of 15(A) of the current limiter is set, and the counted value of the time T is cleared, this routine ending. As a result, only current up to 15(A) comes to be allowed to flow through the DC/DC converter. The procedure from step 707 to step 710 corresponds to the operation occurring from the time td to the time te in the time chart shown in FIG. 6 . [0094] On the other hand, in step 711 to which the routine proceeds if the value of current LiA which flows into the lithium ion battery is determined to be equal to or larger than 15(A) in step 707 or if the determination in step 709 results in T<T 1 , an actual voltage of the lithium ion battery becomes equal to or larger than the control termination voltage V 5 is determined. If the determination in the relevant step determines that the actual voltage of the lithium ion battery becomes equal to or larger than the control termination voltage V 5 , proceed to step 712 , where the current limit value of 15(A) is set on the current limiter, and the counted value of the time T is cleared, this routine ending. [0095] Thus, if the actual voltage of the lithium ion battery is determined to become equal to or larger than the control termination voltage V 5 , even in such a state that the current limit value of the current limiter is released, the current limit of 15(A) is set without any delay. The procedure occurring in step 711 and step 712 corresponds to the operation occurring from the time th to the time t 1 in the time chart shown in FIG. 6 . [0096] In addition, if the actual voltage of the lithium ion battery is determined to be below the control termination voltage V 5 in step 711 , proceed to step 713 . In step 713 , whether or not the target charge voltage TLiV of the lithium ion battery has been changed from the target voltage maximum value V 3 to the normal target voltage V 4 is determined. Then, if the change is determined not to have occurred yet, this routine ends as it is. If the change is determined to have occurred, however, proceed to step 714 , where the current limit value of 15(A) of the current limiter is set, and the counted value of the time T is cleared, the routine ending. [0097] As has been described heretofore, in the fourth embodiment, the value of current which flows into the side of the lithium ion battery via the DC/DC converter while the charge control on the lithium ion battery is carried out is monitored, so that the current limit value of the current limiter of the DC/DC converter is controlled according the current value so monitored. Due to this, compared to the charging method of the related art, the charging accuracy is increased, and the lithium ion battery can be charged completely faster. [0098] Here, the switching of the target charge voltages of the lithium ion battery that has been described in the second and third embodiments or a switching control of the output voltage of the DC/DC converter will be described in greater detail. This switching control corresponds to a control for increasing the value of output voltage of the DC/DC converter and a control for restoring the value of output voltage of the DC/DC converter to its initial value. [0099] FIGS. 8 ( a ), ( b ) are such as to shown a fifth embodiment of the invention and show specifically a first mode of a control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged. FIG. 8 ( a ) is a flowchart which shows the control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged, and FIG. 8 ( b ) is a time chart which shows transitions of the voltage value of the lithium ion battery and the output voltage of the DC/DC converter in the control shown in FIG. 8 ( a ). [0100] In step 801 , whether or not the lithium ion battery has been started to be charged is determined. If the charging of the lithium ion batter has already been started, proceed to step 802 , where the output voltage of the lithium ion battery is switched to the output voltage maximum value V 3 . If the charging of the lithium ion battery has not yet been started, or if the charging thereof has already been executed, proceed from step 801 to step 803 , where whether or not the lithium battery is being charged is determined. Then, if the lithium ion battery is not being charged, the routine ends as it is, whereas if the lithium ion battery is being charged, proceed to step 804 . [0101] An actual voltage LiV of the lithium ion battery is read in step 04 , and a comparison between the actual voltage LiV and the normal target voltage V 4 of the lithium ion battery is carried out to determine whether or not a difference between the normal target voltage V 4 and the actual voltage LiV of the lithium ion battery becomes smaller than a predetermined voltage V 6 in step 805 . The determination made in step 805 results in (V 4 −LiV)>V 6 , the routine ends as it is, whereas if (V 4 −LiV)≦V 6 , proceed to step 806 , where the output voltage of the DC/DC converter is switched from the output voltage maximum value V 3 to the normal output voltage value V 4 and this routine ends. [0102] Part of the procedure where the routine proceeds from step 801 to step 802 in FIG. 8 ( a ) corresponds to a time t 0 in FIG. 8 ( b ), part of the procedure where the routine proceeds from step 804 to step 805 in FIG. 8 ( a ) and the determination in step 805 results in NO corresponds to a waveform drawn from the time t 0 to a time t 1 , and part of the procedure where the routine proceeds from step 805 to step 806 in FIG. 8 ( a ) corresponds to the time t 1 in FIG. 8 ( b ). [0103] FIGS. 9 ( a ) to ( c ) are such as to show a sixth embodiment of the invention and show specifically a second mode of a control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged. FIG. 9 ( a ) is a flowchart which shows the control procedure for the output voltage of the DC/DC converter according to the voltage state of the lead battery when the lithium ion battery is charged, FIG. 9 ( b ) is a time chart which shows a transition of the output voltage of the DC/DC converter when the voltage value of the lead battery is high in the control shown in FIG. 9 ( a ), and FIG. 9 ( c ) is a time chart which shows a transition of the output voltage of the DC/DC converter when the voltage value of the lead battery is low in the control shown in FIG. 9 ( a ). [0104] In the second mode of the control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged, like reference numerals will be imparted to like steps to those of the control procedure of the first mode for description thereof. [0105] In step 801 , whether or not the lithium ion battery has been started to be charged is determined. If the charging of the lithium ion battery has already been started, proceed to step 802 , where the output voltage of the DC/DC converter is switched to the output voltage maximum value V 3 . Then, in step 901 , a voltage value PbV of the lead battery which is resulting then is stored as a voltage value PbVB of the lead battery which results when the lithium ion battery is started to be charged, and this routing ends. [0106] On the contrary, if the charging of the lithium ion battery has not yet been started, or if the charging thereof has already been executed, proceed from step 801 to 803 , where whether or not the lithium ion battery is being charged is determined. Then, if the lithium ion battery is not being charged, the routine ends as it is, whereas if the lithium ion battery is being charged, proceed to step 902 . [0107] In step 902 , an actual voltage PbV of the lead battery is read, and the value so read is stored as a current actual voltage PbVC of the lead battery. In the following step 903 , a comparison is carried out between the voltage value PbVB of the lead battery which resulted when the charging of the lithium ion battery was stared and the current voltage value PbVC of the lead battery to determined whether or not a difference between the voltage value PbVB of the lead battery which resulted when the charging of the lithium ion battery was stared and the current voltage value PbVC of the lead battery becomes larger than a predetermined voltage V 7 . [0108] If the determination made in step 903 results in (PbVB−PbVC)≧V 7 , proceed to step 806 , where the output voltage of the DC/DC converter is switched from the output voltage maximum value V 3 to the normal output value V 4 , and this routine ends. On the contrary, if the determination made in step 903 results in (PbVB−PbVC)<V 7 , proceed to step 904 , where whether or not the current actual voltage PbVC of the lead battery is equal to or smaller than the charge control permissive voltage V 1 for the lead battery is determined. If the determination made in step 904 results in PbVC≦V 1 , proceed to step 806 , where the output voltage of the DC/DC converter is switched from the output voltage maximum value V 3 to the normal output voltage value V 4 , and this routine ends. On the contrary, if the determination made in step 904 results in PbVC>V 1 , the routine ends as it is. [0109] Part of the procedure where the routine proceeds from step 801 to step 802 in FIG. 8 ( a ) corresponds to a time t 0 in each of FIGS. 9 ( b ), ( c ). On the other hand, part of the procedure where the routine proceeds from step 903 to step 806 in FIG. 9 ( a ) corresponds to a time t 1 in FIG. 9 ( b ), and part of the procedure where the routine proceeds from step 903 to step 904 and then to step 806 in FIG. 9 ( a ) corresponds to a time t 1 in FIG. 9 ( c ). [0110] As shown in FIG. 9 ( b ), when the voltage value of the lead battery is high, the output voltage of the DC/DC converter may be switched to the normal output voltage value V 4 at a point in time where the voltage value of the lead battery lowers by the predetermined voltage value V 7 . As shown in FIG. 9 ( c ), however, when the voltage value of the lead battery is low, the voltage value of the lead battery lowers below the charge control permissive voltage V 1 before the voltage value of the lead battery lowers by the predetermined voltage value V 7 . Consequently, in this case, the output voltage of the DC/DC converter needs to be switched to the normal output voltage value V 4 at a point in time where the voltage value of the lead battery lowers below the charge control permissive voltage V 1 . [0111] FIGS. 10 ( a ), ( b ) are such as to show a seventh embodiment of the invention and show specifically a third mode of a control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged. FIG. 10 ( a ) is a flowchart which shows the control procedure for the output voltage of the DC/DC converter according to an integrated value of discharge current of the lead battery when the lithium ion battery is charged, and FIG. 10 ( b ) is a time chart which shows a transition of the integrated value of the discharge current discharged from the lead battery in the control shown in FIG. 10 ( a ). [0112] In the third mode of the control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged, like reference numerals will be imparted to like steps to those of the control procedure of the first mode for description thereof. [0113] In step 801 , whether or not the lithium ion battery has been started to be charged is determined. If the charging of the lithium ion battery has already been started, proceed to step 802 , where the output voltage of the DC/DC converter is switched to the output-voltage maximum value V 3 , and this routine ends. [0114] On the contrary, if the charging of the lithium ion battery has not yet been started, or if the charging thereof has already been executed, proceed from step 801 to step 803 , where whether or not the lithium ion battery is being charged is determined. Then, if the lithium ion battery is not being charged, an integrated value ΣPbA of discharge current of the lead battery is cleared in step 1003 , and this routine ends, whereas if the lithium ion battery is being charged, proceed to step 1001 . [0115] In step 1001 , an integrated value ΣPbA of discharge current of the lead battery is calculated. The integrated value ΣPbA of discharge current of the lead battery is a value which indicates an amount of power that has been discharged from the lead battery since the charging of the lithium ion battery was started. In the following step 1002 , whether or not the integrated value ΣPbA of discharge current of the lead battery becomes equal to or larger than a predetermined threshold value I 2 is determined. [0116] If the determination made in step 1002 results in ΣPbA≧I 2 , proceed to step 806 , where the output voltage of the DC/DC converter is switched from the output voltage maximum value V 3 to the normal output voltage value V 4 , and this routine ends. On the contrary, if the determination made in step 1002 results in ΣPbA<I 2 , this routine ends. [0117] Part of the procedure where the routine proceeds from step 801 to step 802 in FIG. 10 ( a ) corresponds to a time t 0 in FIG. 10 ( b ). On the other hand, part of the procedure where the routine proceeds from step 1002 to step 806 in FIG. 10 ( a ) corresponds to a time t 1 in FIG. 10 ( b ). In this embodiment, in the event that the discharge amount from the lead battery becomes large while the lithium ion battery is being charged, the charging of the lithium ion battery is suppressed so as to suppress the discharge from the lead battery, so that the deterioration of the lead battery is prevented. [0118] FIG. 11 is such as to show an eighth embodiment of the invention and show specifically a fourth mode of a control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged. This embodiment is an embodiment which results by combining together the sixth embodiment and the seventh embodiment, which have been described before, and which is configured such that the output voltage maximum value V 3 of the DC/DC converter is not switched to the normal output voltage value V 4 thereof as long as a voltage drop in the actual voltage PbV of the lead battery does not reach or exceed V 7 and the integrated value ΣPbA of the discharge current of the lead battery does not reach or exceed I 2 after the charging of the lithium ion battery is started a time t 0 . [0119] In the eighth embodiment shown in FIG. 11 , the integrated value ΣPbA of the discharge current of the lead battery comes to reach or exceed I 2 at a time t 1 after the charging of the lithium ion battery was started at a time t 0 , and the voltage drop in the actual voltage PbV of the lead battery comes to reach or exceed V 7 thereafter at a time t 2 , whereupon the output voltage maximum value V 3 of the DC/DC converter is switched to the normal output voltage value V 4 . However, in the event that the voltage drop in the actual voltage PbV of the lead battery reaches or exceeds V 7 earlier at the time t 1 , unless the integrated value ΣPbA of the discharge current of the lead battery comes to reach or exceed I 2 thereafter at the time t 2 , the output voltage maximum value V 3 of the DC/DC converter is not switched to the normal output voltage value V 4 . [0120] The flowchart described with respect to the sixth embodiment and the flow chart described with respect to the seventh embodiment may only have to be combined together to describe a flowchart which shows the control procedure of the eighth embodiment, and hence, in the eighth embodiment, the illustration and description of the relevant flowchart will be omitted. [0121] FIG. 12 is such as to show a ninth embodiment of the invention in which the control procedure of the output voltage of the DC/DC converter when the lithium ion battery is being charged is combined with the control of the current limit value of the current limiter incorporated in the DC/DC converter. In the ninth embodiment, only a flowchart is shown which shows a control procedure. The ninth embodiment is such that in the event that the output voltage of the DC/DC converter is lowered from the output voltage maximum value V 3 to the normal output voltage value V 4 while the lithium ion battery is being charged, at the same time as this occurs, the current limit value of the current limiter incorporated in the DC/DC converter is set to 15(A). The current limit value is such that only one value is set therefor. [0122] In step 1201 , the output voltage of the DC/DC converter is read and stored, and in step 1202 , whether or not the previous output voltage of the DC/DC converter was the output voltage maximum value V 3 is determined. If the previous output voltage of the DC/DC converter was not the output voltage maximum value V 3 , the routine ends as it is, whereas if the previous output voltage of the DC/DC converter was the output voltage maximum value V 3 , proceed to step 1203 . [0123] In step 1203 , whether or not a current output voltage of the DC/DC converter is the normal output voltage value V 4 is determined, and if the current output voltage of the DC/DC converter is not the normal output voltage value V 4 , the routing ends as it is. On the contrary, if the determination made in step 1203 determines that the current output voltage of the DC/DC converter is the normal output voltage value V 4 , proceed to step 1204 while determining that the output voltage of the DC/DC converter has been lowered from the output voltage maximum value V 3 to the normal output voltage value V 4 . In step 1204 , the current limit value of the current limiter incorporated in the DC/DC converter is set to 15(A), and this routine ends. [0124] FIGS. 13 and 14 are such as to show a tenth embodiment of the invention, which shows specifically a fifth mode of a control procedure of the output voltage of the DC/DC converter when the lithium ion battery is charged. FIG. 13 is a flowchart which shows the control procedure for the output voltage of the DC/DC converter according to an integrated value of discharge current of the lead battery when the lithium ion battery is charged and FIG. 14 is a time chart which shows transitions of an integrated value of discharge current discharged from the lead battery, the output voltage of the DC/DC converter and the output voltage value of the lithium ion battery in the control shown in FIG. 13 . [0125] Note that in the fifth mode of the control procedure for the output voltage of the DC/DC converter when the lithium ion battery is charged, like reference numerals will be imparted to like steps to those of the third embodiment for description. [0126] In step 801 , whether or not the charging of the lithium ion battery has already been started is determined. If the charging of the lithium ion battery has already been started, proceed to step 802 , where the output voltage of the DC/DC converter is switched to the output voltage maximum value V 3 , and this routine ends. [0127] On the contrary, if the charging of the lithium ion battery has not yet been started, or when the charging thereof has been executed, proceed from step 801 to step 803 , where whether or not the lithium ion battery is being charged is determined. Then, if the lithium ion battery is not being charged, in step 1003 , an integrated value ΣPbA of discharge current of the lead battery is cleared, and this routine ends, whereas if the lithium ion battery is being charged, proceed to step 1001 . [0128] In step 1001 , an integrated value ΣPbA of discharge current of the lead battery is calculated. In the following step 1002 , whether or not the integrated value ΣPbA of discharge current of the lead battery has reached or exceeded a predetermined threshold value I 2 is determined. If the determination made in step 1002 results in ΣPbA≧I 2 , proceed to step 806 , where the output voltage of the DC/DC converter is switched from the output voltage maximum value V 3 to the normal output voltage value V 4 , and then the routine proceeds to step 1302 . On the contrary, if the determination made in step 1002 results in ΣPbA<I 2 , proceed to step 1301 . [0129] The step 1301 is such as to determine whether or not the charging of the lithium ion battery has been completed, and whether or not the voltage LiV of the lithium ion battery has become the normal output voltage V 4 is determined. Then, if LiV=V 4 , proceed to step 806 , where the output voltage of the DC/DC converter is switched from the output voltage maximum value V 3 to the normal output voltage value V 4 , but if otherwise, this routine ends. [0130] When the output voltage of the DC/DC converter is switched from the output voltage maximum value V 3 to the normal output voltage value V 4 in step 806 , charge current to the lithium ion battery is reduced, and a value of current that is so reduced can be allocated to the charging of the lead battery. Then, in this embodiment, in step 1302 , the charging of the lead battery is carried out, and in step 1303 , an integrated value ΣPbC of charge current to the lead battery is calculated. [0131] Then, in step 1304 , it is determined whether or not the integrated value ΣPbC of charge current to the lead battery has reached I 2 which is the integrated value of the discharge current PbA from the lead battery. If the integrated value ΣPbC of charge current to the lead battery has not reached I 2 , proceed to step 1305 , where it is determined whether or not a switching request has been made which requests the switching of the normal output voltage value V 4 of the DC/DC converter to the output voltage maximum value V 3 thereof. Then, if there has been made no request which requests the switching of the normal output voltage value V 4 of the DC/DC converter to the output voltage maximum value V 3 thereof, the routine ends as it is, whereas even if there has been made a request which requests the switching of the normal output voltage value V 4 of the DC/DC converter to the output voltage maximum value V 3 thereof, in step 1306 , the requested switching is made impermissible, and the routine ends. Namely, in this embodiment, in the event that the integrated value ΣPbC of charge current to the lead battery has not yet reached I 2 , the switching of the output voltage of the DC/DC converter from the normal output voltage value V 4 to the output voltage maximum value V 3 is not permitted. [0132] On the contrary, in the event that the integrated value ΣPbC of charge current to the lead battery has reached I 2 in step 1304 , proceed to step 1307 , and the charging of the lead battery is completed. In this case, in step 1308 , it is determined whether or not there has been made a request which requests the switching of the normal output voltage value V 4 of the DC/DC converter to the output voltage maximum value V 3 thereof, and if there has been made no request which requests the switching of the normal output voltage value V 4 of the DC/DC converter to the output voltage maximum value V 3 thereof, the routine ends as it is, whereas if there has been made a request which requests the switching of the normal output voltage value V 4 of the DC/DC converter to the output voltage maximum value V 3 thereof, the switching is then made permissible in step 1309 , and this routine ends. Namely, in this embodiment, in the event that the integrated value ΣPbC of charge current to the lead battery has reached I 2 , the switching of the normal output voltage value V 4 of the DC/DC converter to the output voltage maximum value V 3 thereof is permitted. [0133] Part of the procedure where the routine proceeds from step 801 to step 802 in FIG. 13 corresponds to a time t 0 in FIG. 14 . In addition, part of the procedure where the routine proceed from step 1002 to step 806 in FIG. 13 corresponds to a time t 1 in FIG. 14 . Furthermore, part of the procedure where the routine proceeds from step 1304 to step 1309 in FIG. 13 corresponds to a time t 2 in FIG. 14 . In addition, part of the procedure where the routine proceeds from step 1301 to step 806 in FIG. 13 corresponds to a time t 3 in FIG. 14 . When the charging of the lithium ion battery is completed at the time t 3 , the lead battery is charged quickly, and the charging of the lead battery is completed at a time t 4 . [0134] In the fifth to tenth embodiments that have been described heretofore, the control is executed in which the output voltages of the DC/DC converter are switched on a timing when the lithium ion battery has been charged to some extent relative to the target charge voltage value, so that the actual voltage of the lithium ion battery is made to converge on the target voltage, whereby not only the lithium ion battery can be charged faster but also the actual voltage of the lithium ion battery can be made to converge on the target voltage in an ensured fashion. In addition, whether or not the charge control is to be executed is not determined by the voltage value of the lithium ion battery, but the state of the lead battery is monitored, and in the event that there occurs a discharge from the lead battery, resulting in the generation of a remarkable voltage reduction on the lead battery side, the charge control made on the lithium ion battery side according to the invention is interrupted, so that the normal charge control is carried out on the lithium ion battery, whereby the deterioration of the lead battery can also be prevented. [0135] In addition, although there occurs no remarkable voltage reduction in the lead battery while the lithium ion battery is being charged, in the event that the voltage value of the lead battery lowers down to the charge control permissive voltage V 1 , since the charge control of the invention is interrupted so as to carry out the normal charge control on the lithium ion battery, the reduction in the voltage of the lead battery to the economical running prohibiting region can be prevented. Furthermore, in the event that a portion of power that was discharged from the lead battery while the charge control of the invention was being executed can be charged back to the lead battery by monitoring the discharge current from the lead battery after the charge control has been shifted from the charge control according to the invention to the normal lithium ion battery charge control, the charge control is restored again from the normal charge control to the charge control according to the invention, thereby making it possible to contribute to the quick charging of the lithium ion battery. [0136] FIG. 15 is such as to show an eleventh embodiment of the invention and is a time chart which explains a correction process in the charge control of the lithium ion battery. FIG. 15 shows a transition of the actual voltage LiV of the lithium ion battery, a transition of the output voltage of the DC/DC converter, a transition of a deviation term (a broken line) between the normal output voltage value V 4 of the target charge voltage value TiLV of the lithium ion battery and the actual voltage LiV of the lithium ion battery and a transition of a value (a solid line) resulting from multiplying the deviation term by a correction coefficient K 1 , a transition of the correction coefficient K 1 , a transition of an integration term of deviation, a transition of a correction coefficient K 2 and a transition of a correction processing cut signal. [0137] When the charging of the lithium ion battery is started at a time t 0 , the actual voltage LiV of the lithium ion battery increases gradually. Since a differential voltage between the actual voltage LiV and the normal target voltage V 4 of the lithium ion battery becomes V 6 at a time t 1 , the charge control is shifted to a control where the actual voltage LiV of the lithium ion battery is made to converge on the target voltage. This control is such that the lithium ion battery is charged so that the actual voltage LiV thereof increases to a certain voltage quickly, and when the lithium ion battery is charged to the vicinity of the target voltage, this control is shifted to another charge control. [0138] Consequently, the output of the DC/DC converter is switched to the normal target voltage V 4 at a time t 1 , and the deviation term and the integration term are used so as to start an output voltage correction process of the DC/DC converter. At the time t 1 , both the value of the correction coefficient K 1 and the value of the correction coefficient K 2 become a predetermined set value, for example, 0.04. Due to this, the output voltage of the DC/DC converter continues to increase slightly until a time t 2 . The correction coefficient K 1 is a coefficient by which the deviation term is multiplied, and the coefficient K 2 is a coefficient by which the integration term is multiplied. The deviation term which is multiplied by the correction coefficient K 1 and the integration term which is multiplied by the correction coefficient K 2 are added to the actual voltage of the lithium ion battery. [0139] Here, consider a state in which the vehicle is accelerated from a time t 2 to a time t 3 . Charging to the lithium ion battery is cut in such a state that the vehicle is accelerated. In this case, in this embodiment, the correction process cut signal is made to be on, and no correction process is executed. In addition, the deviation term and the integration term are cleared so as to be ready for the following correction process. [0140] When the correction process cut signal is made to be off at a time t 3 , the correction process is resumed, and both the value of the correction coefficient K 1 and the correction coefficient K 2 becomes the predetermined set value, for example, 0.04. As a result, the charging to the lithium ion battery is resumed, and the actual voltage LiV of the lithium ion battery increases slightly. Then, in the event that the differential voltage between the actual voltage LiV and the target voltage V 4 of the lithium ion battery becomes equal to or smaller than V 8 at a time t 4 , the value of the correction coefficient K 1 is switched to increase, and the value becomes another correction value, for example, 0.08. This is because in the event that the deviation term is not increased here, it will take a long period of time for the actual voltage LiV of the lithium ion battery to converge on the target voltage V 4 thereof and because the actual voltage LiV of the lithium ion battery is caused to converge on the target voltage V 4 of the lithium ion battery even in the event that the deviation term becomes minute. A transition of the actual voltage LiV of the lithium ion battery is indicated by a broken line which results in the event that the deviation coefficients are not switched over at a time t 4 and the integration term is not cleared at a time t 5 . In this case, the actual voltage LiV of the lithium ion battery overshoots relative to the target voltage V 4 of the lithium ion battery at a certain time after a time t 6 . [0141] In the event that a differential voltage between the actual voltage LiV and the target voltage V 4 of the lithium ion battery becomes equal to or smaller than V 9 thereafter at a time t 5 , the integration term is cleared, and the correction process is carried out by the deviation term only. In the event that the charge control of the invention is carried out over a long period of time, the integration term is increased, and as a result, the output voltage of the DC/DC converter after correction is corrected to the vicinity of the maximum output value, this resulting in the generation of an overshoot in which the actual voltage LiV of the lithium ion battery becomes higher than the target voltage V 4 thereof In order to prevent the generation of this overshoot, when the actual voltage LiV of the lithium ion battery converges on the vicinity of the target voltage V 4 thereof [equal to or larger than (V 4 -V 9 )], the integration term is cleared, and the convergence control is carried out by the deviation term only in which the actual voltage LiV of the lithium ion battery is made to converge on the target voltage V 4 thereof. [0142] By the control that has been described heretofore, the correction process of the output voltage control of the DC/DC converter is executed on the timing when the actual voltage LiV of the lithium ion battery is charged to the certain voltage, and due to the deviation between the actual voltage LiV and the normal target voltage V 4 of the lithium ion battery, by operating the deviation term and the integration term, the accurate and quick charging of the lithium ion battery and the convergence of the actual voltage LiV on the normal target voltage V 4 of the lithium ion battery can be ensured. In addition, even in the event that the deviation term becomes minute, the actual voltage LiV of the lithium ion battery can be made to converge on the normal target voltage V 4 of the lithium ion battery. [0143] Furthermore, when the actual voltage LiV of the lithium ion battery converges on the vicinity of the normal target voltage V 4 of the lithium ion battery, the integration term is cleared, and the correction process is executed by the deviation term only, whereby the likelihood that the actual voltage LiV of the lithium ion battery overshoots the normal target voltage V 4 is reduced. [0144] FIG. 16 is a flowchart diagram illustrating an operation of controlling the cancellation of a limiter value in the process of controlling the charge of the lithium ion battery according to a twelfth example embodiment of the present invention. According to the twelfth example embodiment, the operation is performed as follows: (1) the operation of controlling the early charge of the lithium ion battery is performed only in the case of the voltage of the lithium ion battery being lower and the voltage of the lead battery being higher. (2) The charge of the lead battery has a priority higher than the charge of the lithium ion battery, and when the voltage of the lead battery is low, the operation of controlling the early charge of the lithium ion battery is not performed. (3) When the time period during which the current value flowing through the DC/DC converter is above 15 A, which is an ordinary limiter value, is longer than a predetermined time period, the limiter value is cancelled so that the limiter value is increased to 50 A. [0148] The limiter value is not cancelled at the beginning because, in order to cancel the limiter value, the current may not flow from economical running ECU 10 to the DC/DC converter 3 . Namely, a high level signal must be provided to the terminal DDRE of the DC/DC converter 3 , which will result in unwanted power consumption. In the above twelfth example embodiment, the limiter value is cancelled on the exact timing by detecting the condition which requires that the limiter value to be cancelled, so that the power consumption may be reduced and the lithium ion battery may be charged at an early stage. An example embodiment of this process will be described with reference to FIG. 16 . [0149] In step 1601 , the voltage LiV of the lithium ion battery, the voltage PbV of the lead battery and the motion of a vehicle are sensed. Next, in step 1602 through step 1604 , a drive state of the vehicle is determined. First, in step 1602 , whether the vehicle is in an idle state or not is determined. When the vehicle is in the idle state, step 1605 is performed. Otherwise, whether the vehicle is in a normal drive state or not is determined in step 1603 . When the vehicle is in the normal drive state, step 1605 is performed. Otherwise, whether the vehicle is in a speeding-down drive state or not is determined in step 1604 . If the vehicle is in the speeding-down drive state, step 1605 is performed. Otherwise, step 1618 is performed. [0150] As mentioned above, step 1605 is performed when the vehicle is determined to be in one of the idle state, normal drive state and the speeding-down drive state to determine whether the lithium ion battery is being charge. Here, a signal is inputted to the terminal DDDR of the DC/DC converter of FIG. 1 so that a current flows through the lithium ion battery. When the lithium ion battery is being charged, a step 1606 is performed. In step 1606 , whether a voltage of the lithium ion battery LiV is less than or equal to the voltage V 0 is determined. If the voltage LiV is less than or equal to the voltage V 0 , a step 1607 is performed. In step 1607 , whether a voltage PbV of the lead battery is greater than or equal to the voltage V 1 is determined. [0151] If all the answers to the step 1605 through step 1607 are “Yes”, step 1608 is performed so that the current LiA charged in the lithium ion battery is sensed by a sensor 2 A illustrated in FIG. 1 . If any one of the answers to step 1605 through step 1607 is “No”, step 1618 is performed. When the current LiA of the lithium ion battery is sensed in step 1608 , whether the current LiA of the lithium ion battery is greater than or equal to a reference current I 0 is determined in step 1609 . If the current LiA of the lithium ion battery is greater than or equal to the reference current I 0 , step 1610 is performed. Otherwise, step 1614 is performed. [0152] In step 1610 , a count value of a counter TA that counts the time period during which the current LiA of the lithium ion battery is greater than or equal to the reference current I 0 is increased. In step 1611 , a count value of a counter TB that counts the time period during which the current LiA of the lithium ion battery is smaller than the reference current I 0 is cleared. If the current LiA of the lithium ion battery is greater than or equal to a reference current I 0 , it is determined in step 1612 whether the time period TA during which the current LiA of the lithium ion battery is greater than or equal to a reference current I 0 is longer than or equal to a predetermined time period T 0 . If the time period TA is shorter than the predetermined time period T 0 , the procedure is terminated. If the time period TA is longer than or equal to the predetermined time period T 0 , step 1613 is performed. In step 1613 , the limit value for the current through the limiter installed in the DC/DC converter is cancelled, the procedure is terminated and the control operation of the charge of the lithium ion battery is performed. Accordingly, the limiter value in the DC/DC converter in FIG. 1 is increased from a general current limit value of 15 A to 50 A so that the control operation of the charge of the lithium ion battery is performed. [0153] When the current LiA of the lithium ion battery is lower than the reference current I 0 , step 1614 is performed so that the count value of the counter TA is cleared. Next, in step 1615 , the count value of the counter TB is increased. In the next step 1616 , it is determined whether the time period TB during which the current LiA of the lithium ion battery is less than the reference current I 0 is longer than or equal to a predetermined time period T 1 . If the time period TB is shorter than the predetermined time period T 1 , the procedure is terminated. If the time period TB is longer than or equal to the predetermined time period T 1 , step 1617 is performed. In step 1617 , the limit value for the current through the limiter installed in the DC/DC converter is restored to the ordinary limit value of 15 A so that the control operation of the charge of the lithium ion battery is performed. [0154] In addition, similar to step 1617 , in step 1618 following after one of steps 1604 , 1605 , 1606 and 1607 , the limit value for the current through the limiter of the DC/DC converter is restored to the ordinary limit value of 15 A so that the control operation of the charge of the lithium ion battery is performed. In this case, the count value of the counter TA and counter TB is cleared and terminate the procedure. Based on the control operation described above, with respect to various drive state of a vehicle, the control of the quick charge of the lithium ion battery may be performed when the voltage of the lithium ion battery is lower and the voltage of the lead battery is higher. [0155] FIG. 17 is a flowchart diagram illustrating an operation of converting an output voltage of a DC/DC converter in the process of controlling the charge of the lithium ion battery according to a thirteenth example embodiment of the present invention. According to the thirteenth example embodiment, the operation is performed as follow: (1) the operation of controlling the early charge of the lithium ion battery is performed only in the case of the voltage of the lithium ion battery being lower and the voltage of the lead battery being higher. (2) The charge of the lead battery has a priority higher than the charge of the lithium ion battery, and when the voltage of the lead battery is low, the operation of the controlling the early charge of the lithium ion battery is not performed. (3) During the charging of the lithium ion battery, if excessive discharge of the lead battery does not exist or the voltage of the lithium ion battery does not reach a target charge voltage, the output voltage of the DC/DC converter is set to the maximum voltage V 3 (for example, 16 volts). If the excessive discharge of the lead battery is detected, or the voltage of the lithium ion battery reaches the target charge voltage, the output voltage of the DC/DC converter is returned to the ordinary target voltage V 4 (for example, 15.4 volts). The thirteenth example embodiment will be described with reference to FIG. 17 . The same reference numerals will be used to refer to the same or like steps as those described in the twelfth example embodiment. [0159] In step 1701 , the voltage LiV of the lithium ion battery, the voltage PbV and the current PbA of the lead battery and the motion of a vehicle are sensed. Next, in step 1602 through step 1604 , a drive state of the vehicle is determined similar to the twelfth example embodiment. When it is determined that the vehicle is one of the idle state, normal drive state and the speeding-down drive state, step 1702 is performed. Otherwise, step 1710 is performed. [0160] In step 1702 , which is proceeded when the vehicle is in one of the idle state, normal drive state and the speeding-down drive state, it is determined whether a decision flag VSET that is initialized to “0” has a value of “1” or not. If step 1702 is firstly performed, the decision flag VSET is “0” to proceed to step 1703 . [0161] In step 1703 , the current voltage PbV is stored as a voltage PbVB of the lithium ion battery at an initial stage of charging the lithium ion battery. Next, an accumulation value ΣPbA of the discharge current of the lead battery is cleared and the decision flag VSET described above has a value of “1”. As the decision flag VSET is set to the value of “1”, the answer to step 1702 comes to “Yes” so that step 1703 is not proceeded. Next, in step 1704 , an accumulation value ΣPbA of the discharge current of the lead battery is calculated. [0162] In the next step 1705 through step 1707 , an operation of controlling and converting the target charge voltage of the lithium ion battery to a voltage range between V 4 and V 3 is performed. Here, the voltage V 4 is the ordinary target voltage (output voltage) and the voltage V 3 is the maximum value of the target voltage (maximum output voltage). In step 1705 , a real voltage value of LiV of the lithium ion battery is compared with the ordinary target voltage V 4 to determine whether the difference between the real voltage value of LiV of the lithium ion battery and the ordinary target voltage V 4 is less than a predetermined voltage V 6 . If the difference between the real voltage value of LiV of the lithium ion battery and the ordinary target voltage V 4 (i.e., V 4 −LiV) is greater than or equal to the predetermined voltage V 6 , step 1706 is performed. If the difference between the real voltage value of LiV of the lithium ion battery and the ordinary target voltage V 4 is less than the predetermined voltage V 6 , step 1708 is performed so that the output voltage of the DC/DC converter is converted from the maximum target voltage V 3 to the ordinary target voltage V 4 and the procedure is terminated. [0163] In step 1706 , it is determined whether the accumulation value ΣPbA of the discharge current of the lead battery is greater than or equal to a predetermined threshold value I 2 or not. The accumulation value ΣPbA, which is calculated in step 1704 , of the discharge current of the lead battery represents the amount of power consumption that the lead battery has discharged after the lithium ion battery starts to be charged. If it is determined in step 1706 that the accumulation value ΣPbA of the discharge current of the lead battery is greater than or equal to the predetermined threshold value I 2 , step 1708 is performed to convert the output voltage of the DC/DC converter from the maximum target voltage V 3 to the ordinary target voltage V 4 and the procedure is terminated. If it is determined that the accumulation value ΣPbA of the discharge current of the lead battery is less than the predetermined threshold value I 2 , step 1707 is performed. [0164] In step 1707 , the voltage PbVB of the lithium ion battery at an initial stage of charging the lithium ion battery is compared with a real voltage PbV in the current state of the lead battery to determine whether the difference between the voltage PbVB of the lithium ion battery at an initial stage of charging the lithium ion battery and the real voltage PbV in the current state of the lead battery is greater than or equal to a predetermined voltage V 7 . If the difference (PbVB−PbV) is greater than or equal to the predetermined voltage V 7 , step 1709 is performed to convert the output voltage of the DC/DC converter from the ordinary target voltage V 4 to the maximum target voltage V 3 and the procedure is terminated. If the difference (PbVB−PbV) is less than the predetermined voltage V 7 , step 1708 is performed to convert the output voltage of the DC/DC converter from the maximum target voltage V 3 to the ordinary target voltage V 4 and the procedure is terminated. [0165] In addition, if step 1710 is proceeded following to one of the steps 1604 , 1605 , 1606 and 1607 , the voltage PbVB of the lithium ion battery at an initial stage of charging the lithium ion battery, the accumulation value ΣPbA of the discharge current of the lead battery and the decision flag VSET are all cleared and the procedure is terminated. According to the above control operation, with respect to various drive state, when the voltage of the lithium ion battery is lower and the voltage of the lead battery is higher, a control operation of the charge of the lithium ion battery is performed. Namely, during the charging of the lithium ion battery, if excessive discharge of the lead battery does not exist or the voltage of the lithium ion battery does not reach a target charge voltage, the output voltage of the DC/DC converter is set to the maximum voltage V 3 (for example, 16 volts). If the excessive discharge of the lead battery is detected, or the voltage of the lithium ion battery reaches the target charge voltage, the output voltage of the DC/DC converter is returned to the ordinary target voltage V 4 (for example, 15.4 volts). Therefore, the lithium ion battery may be charged at an early stage. [0166] Note that while in the embodiments, the lead battery and the lithium ion battery are used as the two batteries that are to be installed on the vehicle, the types of batteries used are not limited to these two types, and hence, in place of the lithium ion battery, a nickel-hydrogen battery may be used, or two lead batteries may be used. Furthermore, while in the embodiments, the charge controls are described as being executed by the economical running ECU, the battery ECU and the EFI-ECU which are provided in the automatic engine stop/start control apparatus, the controls described in the invention can also be applied to an integrated ECU in which all controls are executed by an engine control ECU.

A first buttery supplies power to an electrical load of a vehicle when an engine is in a normal operation. A second buttery supplies power to the electrical load of the vehicle when the engine is in an economical operation. A voltage regulator regulates voltages of the first buttery and the second buttery. A current limiter limits a current value flowing from the voltage regulator to the second battery to a limiting value. An engine stopping unit stops the engine when a first predetermined condition is established. An activator starts the engine when a second predetermined condition is established. A current detecting unit detects the current value when a second battery is charged. A changing unit changes a limiting value between a first limiting value and a second limiting value which is higher than the first limiting value according to the detected current value in a case where the voltage of the second battery is no more than a second predetermined value and the voltage of the first battery is no less than a first predetermined value.

PRIORITY This application is a continuation under 35 U.S.C. §120 of U.S. patent application Ser. No. 13/290,931 filed on 7 Nov. 2011, which is a continuation under 35 U.S.C. §120 of U.S. patent application Ser. No. 12/255,998 filed on 22 Oct. 2008, now U.S. Pat. No. 8,054,090, each of which is incorporated herein by reference. BACKGROUND OF THE INVENTION The present invention relates to noise handling in a capacitive touch sensor for detecting proximity of a body, more especially to a capacitive sensor of the so-called active type which is based on measuring the capacitive coupling between a drive and a sense electrode. There are various forms of touch sensitive controls which use a capacitive sensor to sense the presence of a body such as a user's finger. A form of touch sensitive control is disclosed in WO-00/44018. In this example a pair of electrodes are provided which act as a key so that the presence of a body such as a user's finger is detected as a result of a change in an amount of charge which is transferred between the two electrodes. With this arrangement, one of the electrodes (labeled X) is driven with a drive circuit and the other of the pair of electrodes (labeled Y) is connected to a charge measurement circuit which detects an amount of charge present on the Y plate when driven by the X plate. As disclosed in WO-00/44018 several pairs of electrodes can be arranged to form a matrix of sensing areas which can provide an efficient implementation of a touch sensitive two-dimensional position sensor. Such two dimensional capacitive transducing (2DCT) sensors are typically used with devices which include touch sensitive screens or touch sensitive keyboards/keypads which are used in for example in consumer electronic devices and domestic appliances. Devices employing 2DCT sensors have become increasingly popular and common not only in conjunction with personal computers but also in all manner of other appliances such as personal digital assistants (PDAs), point of sale (POS) terminals, electronic information and ticketing kiosks, kitchen appliances and the like. 2DCT sensors are frequently preferred to mechanical switches for a number of reasons. For example, 2DCT sensors require no moving parts and so are less prone to wear than their mechanical counterparts. 2DCT sensors can also be made in relatively small sizes so that correspondingly small, and tightly packed keypad arrays can be provided. Furthermore, 2DCT sensors can be provided beneath an environmentally sealed outer surface/cover panel. This makes their use in wet environments or where there is a danger of dirt or fluids entering a device being controlled attractive. In addition, manufactures often prefer to employ interfaces based on 2DCT sensors in their products because such interfaces are often considered by consumers to be more aesthetically pleasing than conventional mechanical input mechanisms (e.g. push-buttons). Other devices which may incorporate 2DCT sensors include pen-input tablets and encoders used in machinery for feedback control purposes, for example 2DCT sensors are capable of reporting at least a 2-dimensional coordinate, Cartesian or otherwise, related to the location of an object or human body part by means of a capacitance sensing mechanism. Although touch sensitive capacitive sensors such as those described above and disclosed in the above-mentioned disclosures have been successfully deployed in many applications, some applications can present a challenging environment for detecting a change in charge as a result of the presence of a body. For example, the use of a touch sensor on a mobile phone can create a technical problem because there is a variety of disturbing noise signals produced by radio frequency radiation by radio frequency signals and by modulators within the mobile phone. Similarly, a liquid crystal display (LCD) has characteristic switching noise as a result of switching and refreshing pixels. Other types of display may have their own forms of characteristic impulsive noise related to pixel scanning and refresh. Sinusoidal noise, such as that produced by mains electricity may also be present, which can affect the amount of charge detected on a key. This may be significant, for example, when a hand held device such as a mobile telephone is being charged through the mains. FIG. 7 of the accompanying drawings shows an example of sinusoidal noise in the form of a plot of signal strength or amplitude which may be voltage or charge measured with respect to time. Various points 220 are shown to indicate points at which burst measurements are taken for a touch sensor such as those described above. As will be appreciated, as a result of sinusoidal noise represented by a line 222 , an amount of charge transferred from a key by a measurement capacitor of the measurement circuit such as those described above will vary and therefore could in some circumstances cause a false measurement of the presence of a body. FIG. 8 of the accompanying drawings shows another form of noise, namely rectangular or impulsive noise, i.e. noise having high frequency components, such as that which might be produced by switching the pixels in a LCD display with which the touch panel is integrated. A plot of signal strength with respect to time is shown with sampling points 220 , which might be produced by bursts of measurement cycles in a measurement circuit such as those described above. Noise impulses 222 are also shown. If a measurement cycle coincides with a rising edge of a noise impulse, as may arise from an LCD switching event, then an erroneous measurement can be produced which can again cause a touch sensor to erroneously detect the presence of a body. FIG. 9 of the accompanying drawings illustrates this situation showing simultaneous sinusoidal and rectangular noise. As will be appreciated, in general, both sinusoidal noise and rectangular noise may be present during a given time period. Moreover, by its nature, the amount of noise as well as its frequency components will often vary over time. Prior art capacitive sensors adopt a variety of signal processing techniques to filter noise from the acquired signals. For example, boxcar averagers and detection integrators have been used in the past. In principle other types of standard filtering could be used, and may have been used, such as slew rate filters, high frequency pass filters, low frequency pass filters and band pass filters. It is desirable to tailor the filtering depending on how noisy the signals are. For example, it may be desirable to switch certain filters on and off depending on the amount and characteristics of noise, or to set parameters of filters having regard to the noise. SUMMARY OF THE INVENTION The invention provides a method of measuring proximity of a body to a capacitive sensor device comprising a coupling capacitor having a capacitance sensitive to proximity of a body, and a charge accumulation capacitor having first and second plates, the first plate being connected to the coupling capacitor, and the second plate being connected to a voltage output line. A signal measurement is made conventionally by transferring charge indicative of the capacitance of the coupling capacitor from the coupling capacitor to the charge accumulation capacitor, and by reading the voltage on the voltage output line, thereby to measure proximity of a body. An additional noise measurement is made, either before or after the signal measurement, by emulating or mimicking the signal measurement, but inhibiting charge from arising on the coupling capacitor as a result of proximity of a body, thereby to transfer charge indicative of noise induced on the coupling capacitor to the charge accumulation capacitor, and by reading the voltage on the voltage output line, thereby to measure noise. In a preferred embodiment, a method of measuring proximity of a body to a capacitive sensor device is provided, the capacitive sensor device comprising X and Y electrodes forming respective X and Y plates of a coupling capacitor having a capacitance sensitive to proximity of a body, and a charge accumulation capacitor having first and second plates, the first plate being connected to the Y plate of the coupling capacitor, and the second plate being connected to a voltage output line, the method comprising: making a signal measurement by applying one or more cycles of driving the X electrode with an input voltage, thereby to transfer one or more packets of charge indicative of the capacitance of the coupling capacitor from the coupling capacitor to the charge accumulation capacitor, and by reading the voltage on the voltage output line, thereby to measure proximity of a body; and making a noise measurement, either before or after the signal measurement, by emulating or mimicking the signal measurement, but without driving the X electrode, thereby to transfer one or more packets of charge indicative of noise induced on the coupling capacitor to the charge accumulation capacitor, and by reading the voltage on the voltage output line, thereby to measure noise. A noise measurement can therefore be obtained whenever desired, for example after every ‘n’ signal measurements, wherein the noise measurement is taken directly from the coupling capacitors that form the sensing nodes of the capacitive sensor device. The noise measurements should therefore be representative of noise levels when the signal measurements are made, since they are obtained using the same circuit elements and interconnects. Moreover, the noise measurements can be made without any additional hardware requirement on the sensor device. The second plate of the charge accumulation capacitor is preferably connected to a pre-charge line operable to inject a predetermined amount of charge onto the second plate When making the noise measurement, the second plate of the charge accumulation capacitor is then charged before mimicking the signal measurement, thereby to measure noise through its modulation of the amount of charge initially placed on the charge accumulation capacitor. The pre-charge makes the amount of charge accumulated on the charge accumulation capacitor independent of the sign of the charge induced on the Y plate of the coupling capacitor at each cycle of the noise measurement process. This is useful if the controller is not capable of measuring arbitrary voltages from the charge accumulation capacitor. For example, if the controller does not use an analog to digital converter to read off the voltage, but rather a comparator only capable of detecting zero crossings, then the second plate of the charge accumulation capacitor can be pre-charged, e.g. to a certain negative or positive voltage, to ensure that noise contributions will not cause the voltage to change sign before read out, but rather only modulate the pre-charge voltage without changing its sign. In the noise measurement, charge may be transferred to the charge accumulation capacitor in a plurality of cycles, and the voltage on the voltage output line is read after a predetermined number of said cycles. In order that the noise is not averaged out on the charge accumulation capacitor, the number of cycles is preferably 1, or a small integer number such as 2, 3, 4 or 5. In the signal measurement, charge may be transferred to the charge accumulation capacitor in a plurality of cycles, and the voltage on the voltage output line is read after a predetermined number of said cycles. Alternatively, charge may be transferred to the charge accumulation capacitor in a plurality of cycles, wherein the cycles are repeated until the voltage on the voltage output line reaches a threshold value, the number of cycles required to reach the threshold value being read as the measure of the capacitance. In the signal measurement, the number of predetermined cycles could be 1, but is preferably a higher number so that the charge accumulation capacitor has an averaging effect on any noise. For example, the number of cycles is preferably at least 5, 10, 20, 30, 40, 50 or 100. The number of cycles over which charge is transferred from the coupling capacitor to the charge accumulation capacitor is often referred to in the art as the burst length. It will thus be appreciated that the noise measurement will typically use short fixed burst lengths, e.g. of only cycle or a few cycles, whereas the signal measurement will typically use longer burst lengths, which may be fixed, i.e. predetermined in number, or variable, i.e. the number required for the signal to reach a threshold value. Consequently the burst length for the noise measurement will typically be shorter than the burst length for the signal measurement. The method is applied by interspersing signal measurements and noise measurements. Typically, a plurality of signal measurements are made for each noise measurement. The noise measurements may be limited so that they only occupy a certain minority of the total measurement time, for example a noise measurement may be made after ‘n’ signal measurements, where n might be at least 2, 3, 4, 5, 10, 20, 50 or 100. In a 2DCT, a noise measurement might be made once per frame or once every ‘n’ frames, where collection of a frame of touch data is defined as reading a signal measurement from each of the sensing nodes, i.e. each of the coupling capacitors, of the two-dimensional touch panel. The invention also provides a capacitive sensor device comprising a coupling capacitor having a capacitance sensitive to proximity of a body, and a charge accumulation capacitor having first and second plates, the first plate being connected to the coupling capacitor, and the second plate being connected to a voltage output line, the device being operable in a first mode to make a signal measurement by transferring charge indicative of the capacitance of the coupling capacitor from the coupling capacitor to the charge accumulation capacitor, wherein the voltage on the voltage output line is a measure of proximity of a body; and the device being operable in a second mode to make a noise measurement by mimicking the signal measurement, but without enabling charge to arise on the coupling capacitor through proximity of a body, thereby to transfer charge indicative of noise induced on the coupling capacitor to the charge accumulation capacitor, wherein the voltage on the voltage output line is a measure of noise. In a preferred embodiment the invention further provides a capacitive sensor device comprising X and Y electrodes forming respective X and Y plates of a coupling capacitor having a capacitance sensitive to proximity of a body, and a charge accumulation capacitor having first and second plates, the first plate being connected to the Y plate of the coupling capacitor, and the second plate being connected to a voltage output line, the device being operable in a first mode to make a signal measurement by applying successive cycles of driving the X electrode with an input voltage, thereby to transfer successive packets of charge indicative of the capacitance of the coupling capacitor from the coupling capacitor to the charge accumulation capacitor, wherein the voltage on the voltage output line is a measure of proximity of a body; and the device being operable in a second mode to make a noise measurement by mimicking the signal measurement, but without driving the X electrode, thereby to transfer successive packets of charge indicative of noise induced on the coupling capacitor to the charge accumulation capacitor, wherein the voltage on the voltage output line is a measure of noise. The second plate of the charge accumulation capacitor is preferably connected to a pre-charge line operable to inject a predetermined amount of charge onto the second plate, and wherein, in the noise measurement mode, the second plate of the charge accumulation capacitor is pre-charged before mimicking the signal measurement, thereby to measure noise through its modulation of the amount of charge initially placed on the charge accumulation capacitor by the pre-charge. The invention provides benefit to single element touch sensor devices, i.e. devices with one or more isolated buttons, as well as to one-dimensional sensors, such as sliders or scroll wheels, and also to two-dimensional sensors, such as alphanumeric key pad emulators and overlays for displays which thereby provide touch screens. In some embodiments, the device has a plurality of Y electrodes common to each X electrode. In addition, each Y electrode may have its own charge accumulator capacitor. Alternatively, these could be shared through a multiplexer. In a two-dimensional sensor there may be a plurality of X electrodes, and the X and Y electrodes may be arranged to form a two-dimensional array of coupling capacitors distributed over a touch sensitive panel. The device may further comprise a data acquisition unit, such as a microcontroller or other processor, connected to the voltage output line and operable to measure the voltage on the voltage output line to acquire the signal and noise measurements. The device may still further comprise a filter operable to perform numerical processing on the signal measurements, wherein the filter is operable having regard to the noise measurements. The filter may be hosted by the data acquisition unit and provide signal processing in advance of the signal measurements being supplied to higher level systems components. Alternatively, the filter may be hosted in a higher level system component, such as a central processing unit or digital signal processor. The filter preferably has configuration settings that are modified depending on the noise measurements, e.g. a bandpass filter may have its bandpass range changed, a slew filter may have its increment/decrement coefficients changed. Moreover, the noise measurement may govern whether the filter is applied, for example if noise levels are measured to be low, then the filter may be deactivated. BRIEF DESCRIPTION OF THE DRAWINGS For a better understanding of the invention, and to show how the same may be carried into effect, reference is now made by way of example to the accompanying drawings, in which: FIG. 1 a is a schematic block diagram of a touch sensing node, and FIG. 1 b is an example illustration of a user's finger disposed proximate the sensor of FIG. 1 a; FIG. 2 is a schematic block diagram illustrating an electrical equivalent of the touch sensor shown in FIG. 1 b; FIG. 3 is a schematic block diagram of a touch sensing circuit for use with the touch sensing node of FIG. 1 a and FIG. 1 b; FIG. 4 is an example timing diagram illustrating the operation of the sensing circuit shown in FIG. 3 ; FIG. 5 is a circuit diagram illustrating a touch sensitive matrix providing a two-dimensional capacitive transducing sensor arrangement according to an embodiment of the invention; FIG. 6 is a tinning diagram illustrating the operation of the sensing circuit shown in FIG. 3 according an embodiment of the invention; FIG. 7 is a plot of signal strength with respect to time representing a voltage or charge present on a sensing key which has been affected by sinusoidal noise, such as mains noise; FIG. 8 is a plot of signal strength with respect to time representing the voltage or charge on a sensing key in the presence of rectangular noise, such as LCD noise; and FIG. 9 is a plot of signal strength with respect to time representing the voltage or charge on a sensing key in the presence of sinusoidal and rectangular noise. DETAILED DESCRIPTION OF THE INVENTION FIG. 1 a is a schematic cross-section through a touch sensitive control panel 15 in the absence of an actuating body, typically a user's finger or stylus. FIG. 1 b corresponds to FIG. 1 a , but shows the same cross-section in the presence of an actuating body in the form of a user's finger. The touch sensor shown in FIGS. 1 a and 1 b corresponds to an example in which a pair of transverse electrodes form a touch sensor. As shown in FIG. 1 a a pair of electrodes 100 , 104 which form a drive or X plate and a receiving or Y plate in the following description are disposed beneath the surface of a touch sensitive control panel 15 . The electrodes 100 , 104 are disposed beneath a dielectric layer 16 , for example a glass or plastics panel. As shown in FIGS. 1 a and 1 b the touch sensor 10 is arranged to detect the presence of a body such as a user's finger 20 as a result of a change in an amount of charge transferred from the Y plate 104 . As shown in FIG. 1 a when the X plate 100 is charged or driven by a circuit, an electric field is formed which is illustrated by the lines 18 and 19 both above and below the touch panel surface 15 as a result of which charge is transferred to the Y plate 104 . The X plate and the Y plate 100 , 104 form a capacitively chargeable sensing node 10 , referred to as a key in the following. As shown in FIG. 1 b as a result of the disturbance of the electric field 18 due to the presence of the user's finger 20 the electric field above the surface of the control panel 15 is disturbed as a result of an earthing or grounding effect provided by the user's finger 20 as illustrated schematically by ground 34 . An equivalent circuit diagram of the touch sensor shown in FIGS. 1 a and 1 b is shown in FIG. 2 . In FIG. 2 equivalent capacitances are illustrated in the form of a circuit diagram. A capacitance formed between the X plate 100 and the Y plate 104 of the key is a capacitance CE 105 (sometimes also referred to as Cx in the art) which is in effect a coupling capacitor. The presence of the body 20 has an effect of introducing shunting capacitances 30 , 32 , 33 which are then grounded via the body 20 by an equivalent grounding capacitor 22 to the ground 34 . Thus the presence of the body 20 affects the amount of charge transferred from the Y plate of the key and therefore provides a way of detecting the presence of the body 20 . This is because the capacitive coupling between the X plate 100 and the Y plate 104 of the key CE 105 reduces as a result of the T-bridge effect caused by the increased capacitance 33 . It will be appreciated by the skilled person that FIGS. 1 a and 1 b are depicting a so-called active capacitive sensors based on measuring the capacitive coupling between two electrodes (rather than between a single sensing electrode and a system ground). The principles underlying active capacitive sensing techniques are described in U.S. Pat. No. 6,452,514. In an active-type sensor, one electrode, the so called drive electrode, is supplied with an oscillating drive signal. The degree of capacitive coupling of the drive signal to the sense electrode is determined by measuring the amount of charge transferred to the sense electrode by the oscillating drive signal. The amount of charge transferred, i.e. the strength of the signal seen at the sense electrode, is a measure of the capacitive coupling between the electrodes. When there is no pointing object near to the electrodes, the measured signal on the sense electrode has a background or quiescent value. However, when a pointing object, e.g. a user's finger, approaches the electrodes (or more particularly approaches near to the region separating the electrodes), the pointing object acts as a virtual ground and sinks some of the drive signal (charge) from the drive electrode. This acts to reduce the strength of the component of the drive signal coupled to the sense electrode. Thus a decrease in measured signal on the sense electrode is taken to indicate the presence of a pointing object. FIG. 3 provides a circuit diagram, which forms a touch sensor by sensing an amount of charge transferred from the X plate 100 shown in FIG. 2 to the Y plate 104 and includes a charge measurement circuit which has been reproduced from WO-00/44018, which corresponds to U.S. Pat. No. 6,452,514. As shown a drive circuit 101 is connected to the X plate 100 of the key 105 and the Y plate 104 of the key 105 is connected to an input 106 of a charge measurement circuit 108 , wherein the X and Y plates collectively form the capacitor 105 . The input 106 is connected to a first controllable switch 110 and to one side of a measuring capacitor Cs 112 on which charge is accumulated as a measure of capacitive coupling. The other side of the measurement capacitor 112 is connected via a second controllable switch 114 to an output 116 of the measurement circuit 108 which is fed as a voltage VOUT to a controller 118 . A first input control channel 103 is used to control the operation of the drive circuit 101 . The first and second controllable switches 110 and 114 are controlled by the controller 118 through respective first and second switch control lines 146 and 148 Similarly, the drive circuit 101 is controlled by the controller 118 through the first input control channel 103 . In the illustrated circuit diagram, a convention has been adopted whereby a control input of each of the switches 110 , 114 is open when the control input is “0” and closed when the control input is “1”. The other side of each of the switches 110 , 114 is connected to ground, so that if the control input is “1” then the connecting input would be connected to ground. A similar convention has been adopted for drive circuit 101 , whereby when the control input 103 is “0” the X plate is connected to ground and when the control input is “1” the X plate is connected to a reference voltage “V”. FIG. 4 is a timing diagram which shows the operation of the touch sensor, and in particular the function of the measurement circuit arranged to measure the amount of charge transferred from the X plate 100 to the Y plate 104 of the key 105 . Four timing diagrams 130 , 132 , 134 , 138 are shown to illustrate the operation of the measurement circuit 108 . A first timing diagram 130 represents the control input applied to the second switch 114 . On the left hand side, the logical value of the control input is shown, whereas on the right hand side the effect at the connecting point 114 . 1 is shown to be either “Z” in which the connecting point 114 . 1 is isolated or floating, or for a logical control input of 1 grounded. A second timing diagram 132 represents the control input applied to the first switch 110 . The logical control input values “0” or “1” of a connecting point 110 . 1 are shown at either floating (Z) or ground (0). A third timing diagram 134 shows a relative timing of a drive signal provided to the X plate 100 of the key in which case, in contrast to the timing diagrams 130 , 132 for the two switches 110 , 114 , the value of the timing diagram is an absolute value so that the left hand side illustrates that the voltage varies between ground and the reference voltage “V”, which is the voltage used to charge the X plate 100 . A fourth timing diagram 138 provides an illustration of the example signal strength or voltage produced on the measurement capacitor 112 as a result of the opening and closing of the switches 110 , 114 and the driving of the X plate 100 in accordance with the timing illustrated by the timing diagrams 130 , 132 , 134 . The timing diagrams 130 , 132 , 134 , 138 will now be explained as follows: At a first time t 1 , the charge measurement circuit 108 is initialised, i.e. reset, using control lines 146 and 148 for respective switches 110 and 114 being high (1) and control line 103 for drive circuit 101 being low (0). Control lines 146 , 148 , 103 are lines connected to the controller 118 . The Y plate 104 , the X plate 100 and the charge measurement capacitor 112 are thus set to ground. Correspondingly, the output voltage across the charge measurement circuit 112 is at zero. It will be appreciated that connections to ground and VDD could be reversed in other embodiments. At a second time t 2 the logical input to the second switch 114 is set low (0), thereby opening the switch and floating the connecting point 114 . 1 . At a third time t 3 the control input to the switch 110 is set low (0), thereby floating the connecting point 110 . 1 , which is YA before, at a time t 4 the control input 103 of the drive circuit 101 is set high (1), thereby the X plate 100 of the key 105 is connected to the reference voltage “V”. Then, in order to charge the measurement capacitor Cs for a period S between t 5 and t 6 , the control input to the second switch 114 is set high (1), thereby grounding YB to transfer charge induced on the Y plate of the key 105 onto the charge measurement capacitor 112 , until time t 6 , when the control input to the second switch 114 is set low (0), which again floats the connecting point 114 . 1 . After charging the measurement capacitor Cs for a first dwell time between t 5 and t 6 , at t 7 the control input to the first switch 110 is set high (1), thereby grounding the connecting point 110 . 1 , which is connected to the other side of the charge measurement capacitor Cs 112 . As a result, the voltage across the measurement capacitor can be measured. The amount of charge transferred from the Y plate 104 onto the measurement capacitor Cs 112 during the dwell time between t 5 and t 6 is represented as the output voltage VOUT. At time t 8 the control input 103 to the drive circuit 101 goes low (0), thereby the X plate 100 of the key 105 is connected to ground which concludes a first measurement cycle. At time t 9 the next measurement cycle of a measurement burst occurs. At t 9 the control input to the switch 110 goes low (0) thereby floating YA, before the control input 103 to the drive circuit 101 again goes high (1), thereby connecting the X plate 100 to the reference voltage “V”, at time t 10 . The measurement capacitor 112 is again charged from charge transferred from the Y plate 104 of the key onto the measurement capacitor 112 . As with the first burst at point t 11 the control input to the switch 114 goes high (1) thereby grounding the point 114 . 1 and driving charge onto the measurement capacitor until t 12 , when the control input to the switch 114 goes low, again floating YB. Thus again charge is transferred from the Y plate 104 during the dwell period between t 11 and 02 , thereby increasing the voltage across the measurement capacitor Cs as represented as the output voltage VOUT. It will be appreciated that in FIG. 4 , VOUT is shown with respect to a ground connection at YA, i.e. if point 110 . 1 is ground. At t 13 the control input to the switch 110 is set high (1) thereby grounding YA and at t 14 control input 103 to the drive circuit 101 goes low (0), thereby connecting the X plate 100 of the key 105 to ground which concludes the second measurement cycle. Thus, as with the first measurement cycle, an amount of charge has been transferred from the Y plate, which has then increased the voltage across the measurement capacitor 112 , which represents an amount of charge transferred from the Y plate. After ‘n’ measurement cycles of a burst, ‘n’ packets of charge will have been transferred from the Y plate to the measurement capacitor 112 where the charge is accumulated. By bursting in this way the signal from each transfer is averaged on the measurement capacitor 112 , so that when it is read out noise reduction has effectively already taken place, as is well known in the art. The amount of charge on the measurement capacitor 112 is determined with the aid of a resistor 140 . One side of the resistor 140 is connected to the measurement capacitor 112 and the other side, labeled SMP, is connected to a controllable discharge switch 142 . The discharge switch 142 is connected to receive a control signal from the controller 118 via a control channel 144 . The discharge switch 142 floats SMP when in position ‘0’ and connects SMP to a voltage VDD when in position ‘1’. The float position is selected during measurement, i.e. during the cycles illustrated in FIG. 4 , and the VDD position is selected to discharge the measurement capacitor Cs 112 through the discharge resistor 140 after charge accumulation through a number of cycles. The controller 118 is operable to determine the amount of charge accumulated on the measurement capacitor by measuring the amount of time, e.g. by counting the number of clock periods, it takes for the charge on the measurement capacitor Cs to discharge, i.e. the amount of time for the voltage VOUT to reduce to zero. The number of clock periods can therefore be used to provide a relative signal sample value for the respective measured charge signal. This is preferably achieved using a comparator that compares VOUT to a zero or ground signal, while capacitor Cs is discharged. Using the same principles of construction and operation, a matrix or grid of touch sensitive switches can be formed to provide a 2DCT sensor. A user can then at a given time select one, or in some cases a plurality of, positions on a sensitive area. FIG. 5 shows a two-dimensional touch sensor employing an array of sensors of the kind described with respect to FIGS. 3 and 4 . The charge measurement circuit is used in conjunction with drive circuits 101 . 1 , 101 . 2 , 101 . 3 , 101 . 4 and is arranged to drive different sensor points 205 . As shown each of the drive circuits 101 . 1 , 101 . 2 , 101 . 3 , 101 . 4 is controlled by the controller 118 to drive each of the corresponding lines X 1 , X 2 , X 3 , X 4 , using first control inputs 103 . 1 , 103 . 2 , 103 . 3 , 103 . 4 in the same way as the X plate 100 is driven in FIG. 3 and represented in FIG. 4 . Similarly, an input 107 provides a reference voltage “V”. The output of the coupling capacitors at each of the points 205 are connected to one side of measuring capacitors Cs 112 . 1 , 112 . 2 , 112 . 3 , 112 . 4 which are arranged to measure an amount of charge present on the Y plate, Y 1 , Y 2 , Y 3 , Y 4 providing output signals 116 . 1 , 116 . 2 , 116 . 3 , 116 . 4 to detect the presence of an object in the same way as the operation of the circuit shown in FIG. 3 and FIG. 4 . This is achieved by applying control signals to the switches 110 a , 110 b , 110 c , 110 d , 114 a , 114 b , 114 c , 114 d in a corresponding manner to the arrangement explained above with reference to FIGS. 3 and 4 . In FIG. 5 , some of the detail from FIG. 3 has been omitted for clarity. In particular the resistor 140 , its switch 142 and actuating line 144 are not shown. It is noted that each of the switches 142 can be commonly actuated by a single actuating line 144 from the controller 118 , since they only need to be switched together to perform their function described above. More details for the operation of such a matrix circuit are disclosed in WO-00/44018. The controller operates as explained above to detect the presence of an object above one of the matrix of keys 205 , from a change in the capacitance of the keys, through a change in an amount of charge induced on the key during a burst of measurement cycles. However, the presence of a noise signal can induce charge onto the keys of a touch sensor and provide a false detection or prevent a detection being made. FIG. 6 shows timing diagrams 130 , 132 , 134 , 138 and 600 used to illustrate the further operation of the of the measurement circuit 108 shown in FIG. 3 to reduce the effects of noise according to an embodiment of the invention. The tinning diagrams 130 , 132 , 134 and 138 correspond to the same timing diagrams shown on FIG. 4 , except that the timing diagrams shown in FIG. 6 precede those shown in FIG. 4 , i.e. they run from t- 10 to t 1 . The same convention for each of the switches 110 and 114 shown in FIGS. 3 and 4 is adopted here. The convention for the switch 142 , corresponding to the timing diagram 600 , is also shown in FIG. 3 . The control signal 144 connects the SMP resistor 140 to either float (Z) or VDD depending on its position which is determined by the control line 144 . When the control line 144 is high (1) switch 142 connects SMP resistor 140 to the VDD. When the control line is low (0) switch 142 connects resistor 140 to float (Z). The timing diagrams 130 , 132 , 134 and 138 for YB, YA, X and VOUT respectively correspond to those shown in FIG. 4 and described above. Timing diagram 600 represents the control input applied to the switch 142 . On the left hand side, the logical value of the control input is shown, whereas on the right hand side the effect at the SMP resistor 140 is shown to be either “VDD” in which the resistor 142 is connected to VDD or float (Z) in which the resistor 142 is floated. The timing diagrams 130 , 132 , 134 , 138 and 600 will now be described. In FIG. 6 at a first time t- 10 , the charge measurement circuit 108 is initialized with both the control inputs for the switches 110 , 114 being high (1) so that both the Y plate and the charge measurement capacitor 112 are set to ground and the input to the switch 142 is set low (0) so that the resistor 140 is floated. This is to discharge the capacitor. At time t- 9 , the logical input to the control switch 114 is set to zero, thereby opening the switch and floating the connecting point 114 . 1 which connects the output voltage 116 to one side of the measurement capacitor Cs 112 . At a next time t- 8 the control input to the switch 142 is set to high (1), thereby connecting the resistor 140 to VDD. The measurement capacitor Cs 112 is now connected to ground at one side and to VDD via the SMP resistor 140 on the other side, thereby charging the capacitor. The measurement capacitor Cs 112 is charged for a predetermined time T to charge Cs by a set amount. This direct charging of the measurement capacitor Cs bypassing the X drive circuit 101 is referred to as a pre-charging. The pre-charge on the capacitor can be calculated, since the voltage, the value of the capacitor Cs 112 , the value of the resistor 140 and the charge time T are all known. VOUT 138 shows the voltage of the capacitor between t- 8 and t- 7 . At time t- 7 the control input to the switch 142 is set low (0) thereby floating the SMP resistor 140 . At time t- 6 the control input to switch 110 is set low (0) thereby point 110 . 1 is floating. To measure the amount of noise on the coupling capacitor, switches 110 , 114 and 142 are driven in the same way as if they were collecting signal measurements, i.e. the same way as described with reference to FIG. 4 , with the exception that the X drive circuit 101 is not driven. Namely, as can be seen from FIG. 6 , the control switch to the X plate remains low (0) from time t- 10 to t 1 , so no voltage is applied to the X plate 100 of the coupling capacitor 105 . As a result, the charge accumulated on the measurement capacitor 112 is the pre-charge modulated by packets of charge picked up on the Y plate 104 of the coupling capacitor 105 during the acquisition cycles. This modulation will be a measure of noise, since it can only have been noise that has induced charge on the coupling capacitor 105 , given the absence of X drive signal during the dwell times of each cycle. To summarize, a noise measurement is made by mimicking the signal measurement, but without driving the X plate 100 . Moreover, by pre-charging the accumulation capacitor before mimicking a normal signal measurement without driving the X plate, noise is measured through its modulation of the amount of charge initially placed on the charge accumulation capacitor by the pre-charge. In order to charge or discharge the measurement capacitor Cs for a period S between t- 5 and t- 4 , the control input to the switch 114 is set high (1) thereby grounding YB to transfer charge induced on the Y plate 104 of the key 105 , due to noise, onto the charge measurement capacitor 112 , until t- 4 when the control input to the switch 114 is set to low (0), which again floats the connecting point 114 . 1 . During time t- 5 to t- 4 noise is accumulated on the capacitor Cs as shown on VOUT 138 . The noise accumulated on the capacitor Cs during the dwell time could result in the voltage at time t- 4 being higher or lower than the voltage at time t- 5 . Two different outcomes are illustrated in FIG. 6 for VOUT 138 at time t- 4 . These are illustrated as solid line 602 , between tunes t- 5 and t- 1 and dotted line 604 between time t- 5 and t- 1 . The solid line 602 illustrates a noise signal that has removed charge from the measurement capacitor Cs and the dotted line 604 illustrates a noise signal that has added charge to the measurement capacitor Cs. After charging the measurement capacitor Cs for a first dwell time S between t- 5 and t- 4 , at t- 3 the control input to switch 110 is set high (1), thereby grounding the connecting point 110 . 1 (YA), which is connected to the other side of the charge measurement capacitor Cs 112 . This will allow the charge on the capacitor Cs to be measured. The measurement of the charge accumulated on capacitor Cs is carried out in the same manner as described above for measuring the charge accumulated on the capacitor due to a touch. At time t- 2 the control input 144 to switch 142 goes high (1), thereby connecting the SMP resistor to VDD. As a result, the voltage across the measurement capacitor Cs 112 can be measured between times t- 2 and t- 1 . The amount of charge transferred from the Y plate 104 onto the measurement capacitor Cs 112 during the dwell time between t- 5 and t- 4 in addition to the pre-charged charge on the measurement capacitor is measured using the method described above. The read-out time U is used to determine the amount of charge on the capacitor Cs, in addition to the capacitance of the capacitor Cs, which is known. It will be appreciated that the read-out time for each of the two outcomes represented by the solid line 602 and the dotted line 604 will be different i.e. the time to discharge the measurement capacitor Cs for the scenario illustrated by the dotted line 604 will be higher than for the scenario represented by the solid line 602 . However, for simplicity these have been shown having the same read-out time in FIG. 6 . At time t 0 the control switch 114 is set high (1) so that the timing diagram continues on to time t 1 as shown in FIG. 3 . The control input for switch 142 is held low (0) during the measurement of the charge on the Y plate while the X plate is driven, as described above. The process above is only described for a single charge measurement circuit. However, it will be appreciated that the process described above could be carried out on each of the Y plates using each of the charge measurement circuits connected to the Y plates as shown in FIG. 5 . In the above example, the noise acquisition precedes the signal acquisition. However, the time sequence is arbitrary. In practice, the measurements will be interspersed, with for example one noise sample being taken followed by several signal samples and then another noise sample etc. Once the charge on the measurement capacitor Cs due to noise on the Y plate is measured, the amount of noise is determined. The detected noise is the difference between the charge on the measurement capacitor from the pre-charge cycle and the measured charge on the measurement capacitor after the dwell time. It will be appreciated that the charge on the capacitor after the dwell time could be less than or equal to the pre-charge charge, since charge can be removed from the capacitor as a result of the noise as well as being added. In other words the noise is a charge value that is obtained from subtracting the amount of charge on the measurement capacitor Cs 112 after the pre-charge step from the measured charge on the measurement capacitor Cs 112 after the dwell time. “Noise charge” will be used to identify this charge difference. A noise factor is now calculated using the calculated noise charge. To calculate the noise factor, the detected noise charge from the current sample and the previous four samples are used. The standard deviation of these 5 samples is calculated. This will be referred to as the noise factor. It will be appreciated that it is computationally intensive to calculate a square root, so in the preferred embodiment the square of the standard deviation is used. It will be appreciated that more or less samples could be used to obtain the noise factor. Furthermore, it will be appreciated that other methods method for averaging the noise charge could be used. The method described above has shown how a noise factor can be calculated using the touch sensor and touch matrix shown in FIGS. 3 and 5 . Examples of how the noise factor is used to reduce the effects of noise in such an arrangement shown in FIGS. 3 and 5 are now described. In an embodiment of the invention the controller 118 contains a single filter. The filter could be any type of linear or non-linear filter, for example a low pass filter. The noise factor is used to control whether or not the filter is used. If the noise factor is below a specified threshold the filter is not used. However, if the noise factor is above a specified threshold the filter is applied to the measured signals. For example, if the touch matrix was implemented in a mobile phone, the noise factor during normal operation may fall below the specified value, thereby no filtering is used. However, if the mobile phone is connected to a phone charger during operation, for example, the phone charger may introduce noise. Therefore, if the noise factor due to the phone charger exceeds the specified value the filter would be applied to the measured signal. Once the phone battery is charged and the phone is disconnected from the charger the noise factor will again be below the specified value and the filter will not be applied to the measured signal The embodiment of the invention provides a method for configuring one or more filters in response to real time noise signals present on the touch matrix. The sensing hardware shown in FIGS. 3 and 5 typically remains unchanged. Therefore this may be implemented as a firmware update. Since, the method described above provides a method for activating or deactivating filters, when there is no noise source or the noise source is low i.e. below a specified value, the time taken to determine the location of a touch on a touch matrix can be increased, since there is no requirement for filtering. In another embodiment the controller 118 may contain a low pass filter that can be configured using the acquired noise factor using the method described above. The frequency pass of the filter may be adjusted based on the amount of noise detected. For example, if the expected signal level detected on the touch matrix is ‘S’, the frequency pass of the filter during normal operation could be ‘S’+10. However, if the detected noise signal is very high for example, then the frequency pass of the filter band could be increased to accommodate a signal with a greater amount of noise. Alternatively, the low pass filter could be replaced with a slew rate limiter. The rate at which the input value is allowed to change is adapted in accordance with the noise factor. If the noise factor is high, the slew rate limiter is adapted to allow signals that change more slowly, i.e. the allowed slew rate is low. Alternatively, if the noise factor is low, the slew rate limiter is adapted to allow signals that change more quickly, i.e. the allow slew rate is high. It will be appreciated that other types of slew rate limiter could be used other than a linear changing slew filter. For example where two consecutive samples exceed a predetermined value, the rate of change can be capped to a fixed increment or decrement so that the slew rate limiter can settle on the average value more quickly. In summary of the above-described embodiments, a capacitive touch sensor is provided for detecting proximity and location of a body, the sensor comprising: one or multiple X lines; a plurality of Y lines each arranged to have a portion thereof adjacent to a portion of each of the X lines to form a plurality of sensing capacitors; a charge measurement capacitor connected to each Y line; one or more drive circuits arranged to drive respective ones of the X lines to enable charge transfer from the sensing capacitors associated with that X line to the respective measurement capacitors connected to the Y lines. The touch sensor is operable under control of a controller to measure a signal value in the usual way from one of the measurement capacitors, the signal value being indicative of the amount of charge transferred from one of the plurality of Y lines when actuating one of the respective X lines. The touch sensor is further operable under control of a controller to measure a noise value from one of the measurement capacitors indicative of the amount of charge transferred from one of the plurality of Y lines without first actuating one of the respective X lines. Charge is preferably injected onto the charge measurement capacitor before measuring the noise value, so that any noise contribution is accurately measured.

In particular embodiments, an apparatus includes a charge-measurement capacitor having a first plate coupled to a second plate of a coupling capacitor and a non-transitory computer-readable storage medium embodying logic that is operable when executed to ground a first plate of the coupling capacitor; inject a pre-determined amount of charge onto the charge-measurement capacitor; and transfer an amount of charge accumulated on the second plate of the coupling capacitor to the first plate of the charge-measurement capacitor. The charge accumulated on the second plate of the coupling capacitor is due at least in part to noise. The logic is also operable when executed to determine, through a measured voltage across the charge-measurement capacitor, the amount of charge.

BACKGROUND OF THE INVENTION The present invention relates to a material having a low dielectric constant (low k material) useful as an insulating film used for interlayer insulation of semiconductor elements, as a barrier metal layer or an etch stopper layer, or as a substrate for electric circuit parts, and also relates to an insulating film comprising this material and a semiconductor device having the insulating film. Demands for high integration and high speed of semiconductor devices are increasing more and more. In order to meet these demands, there have been made a study on conductive layer materials having a lower electric resistance than conventional aluminum alloy, namely a study on wiring materials, and a study on insulating layer materials having a lower dielectric constant than conventional silicon oxide. In particular, these materials are needed in wiring of semiconductor devices if the structural minimum dimension of the semiconductor devices becomes smaller than about 0.18 μm, as known, for example, from “Recent Development in Cu Wiring Technology” edited by S. Shinmiyahara, N. Awaya, K. Ueno and N. Misawa published by Realize Company, Japan in 1998. FIG. 5 is a section view showing a two layer copper wiring structure in a semiconductor device disclosed in the above publication. In the figure, numeral 1 is a silicon substrate, and numeral 2 is a first insulating layer having trench 3 corresponding to a first wiring pattern. The first insulating layer 2 is made of a silicon oxide film having a dielectric constant of 4.2 or a fluorine-containing silicon oxide film having a dielectric constant of 3.2 to 3.5. Further, studies have been made on applicability, as alternates, of materials having a lower dielectric constant than 2.8 such as silicon-based inorganic polymer materials, organic polymer materials, amorphous fluorine-containing carbon films and porous silicon oxide films. The bottom and the side faces of trench 3 are covered with first conductive film 4 having a diffusion preventive function as a barrier metal. As the first conductive film 4 is used titanium nitride (TiN), tantalum nitride (TaN), tungsten nitride (WN), or a trinary barrier metal comprising each of these nitrides and silicon. First copper conductive layer 5 is formed to fill the trench 3 covered with the first conductive film 4 . Numeral 6 is a first insulating film having a diffusion preventive function against copper, which is made of silicon nitride. Numeral 7 is a second insulating layer, which is made of a material similar to that of the first insulating layer 2 . A hole 8 is formed in the first insulating film 6 and the second insulating layer 7 therethrough, and the bottom and side surfaces of the hole 8 are covered with a second conductive film 9 having a diffusion preventive function and contacting the first copper conductive layer 5 . The hole 8 which is covered with the second conductive film 9 is filled with a second copper conductive layer 10 . A trench 12 corresponding to a second wiring pattern is also formed in the second insulating layer 7 , and the inner surfaces of trench 12 are covered with third conductive film 11 having a diffusion preventive function. The trench 12 which is covered with the third conductive film 11 is filled with a third copper conductive layer 13 . The second and third conductive films 9 and 11 are made of a material similar to that of the first conductive film 4 . The upper surface of the third copper conductive layer 13 is covered with a second insulating film 14 made of silicon nitride having a diffusion preventive function against copper. The first and third copper conductive layers 5 and 13 constitute wiring in the lower layer and wiring in the upper layer respectively, and the second copper conductive layer 10 electrically connects these wirings in the upper and lower layers therebetween. While the wiring of two layer structure is shown in FIG. 4 , this structure may be repeatedly stacked to form a multi-layer structure. The wiring structure shown in FIG. 5 is formed through a so-called Damascene process, which will be described below. Trench 3 corresponding to a wiring pattern is formed in first insulating layer 2 , and first conductive film 4 which serves as a barrier metal, is formed on the inner surface of the trench 3 . A copper film is is then formed on the first insulating layer 2 to fill the trench 3 . Unnecessary barrier metal and copper films formed on portions other than the trench 3 are removed by CMP (chemical mechanical polishing) to leave the barrier metal and copper only in the trench 3 to form first copper conductive layer 5 . In such a manner, the copper wiring in the lower layer is formed in the trench 3 with the bottom and side surfaces thereof covered with the first conductive film 4 . Then, silicon nitride film 6 and second insulating layer 7 are sequentially stacked on the first insulating layer 2 . Trench 12 having a pattern corresponding to the second wiring and hole 8 extending to the first copper conductive layer 5 are formed in the silicon nitride film 6 and the second insulating layer 7 therethrough. Second and third conductive films 9 and 11 are formed as the barrier metal on the surfaces of the trench 12 and the hole 8 . The trench 12 and the hole 8 are then filled with copper by copper film forming, followed by removal of unnecessary copper and barrier metal on the second insulating layer 7 using CMP to thereby form the wiring in the upper layer. Thereafter, second insulating film 14 is formed. In case that a polymeric material or a porous silicon oxide, which have a lower dielectric constant than silicon oxide and fluorine-containing silicon oxide, is used as a material for the first or second insulating layer or the first or second insulating film of semiconductor devices having the above wiring structure, a problem arises about deterioration in reliability of wiring and device, since these materials have a lower thermal conductivity as compared with conventionally used silicon oxide and heat generation in a wiring may cause temperature rise of semiconductor devices. FIG. 6 is a section view showing a wiring structure in a semiconductor device disclosed in W. Y. Shih, M. C. Chang, R. H. Havemann and J. Levine, Symposium on VLSI Technology Digest, pages 83-84, 1997, wherein two kinds of insulative materials are used in the above-mentioned first and second insulating layers respectively in order to solve the problem associated with poor thermal conductivity. That is to say, a material having a low dielectric constant such as a polymeric material is used as a material of insulating layers 15 and 16 in which wiring is formed by each of first copper conductive layer 5 and third copper conductive layer 13 . On the other hand, silicon oxide which has a good thermal conductivity and has been conventionally used as an insulating material of a wiring-forming layer, is used as a material of insulating layer 18 in which hole 8 is formed and as a material of insulating layer 17 disposed between first copper wiring 5 and substrate 1 , thereby suppressing deterioration in thermal conductivity as a whole. Numerals 4 , 9 and 11 denote first, second and third conductive films respectively which are formed as a barrier metal. Numeral 10 is a second copper conductive layer filled in the hole 8 . Numeral 12 is a trench, and numeral 14 is a second insulating film. In the former publication, it is described that owing to scale down of pattern size associated with high integration of integrated circuits in semiconductor devices and owing to increase in wiring length resulting from increase in chip area, propagation delay of signals on wiring is growing to a major cause hindering advent of high speed devices. Solution of such a problem would require reduction in wiring resistance and use of insulating films having low dielectric constant for reduction in electrostatic capacitance between wirings, namely reduction in wiring capacitance. As a wiring material for this purpose, copper is beginning to be used in place of aluminum alloy used conventionally. On the other hand, as an interlayer dielectric for this purpose, a fluorine-containing silicon oxide film having a dielectric constant of 3.2 to 3.5, namely SiOF, is also beginning to be used in place of silicon oxide having a dielectric constant of 4.2. However, in case of forming an interlayer insulating film from SiOF, its dielectric constant is from about 3.2 to about 3.5 and, therefore, the reduction in capacity between wirings and the prevention of propagation delay of signals on wiring are not sufficiently achieved, although the dielectric constant of interlayer insulating film becomes lower than conventional one. With respect to interlayer insulating films formed from organic compounds, dielectric constant of 2.7 is achieved by a film of a polyimide into which fluorine atom is introduced or by an aryl ether polymer, but they are still unsatisfactory for use as an interlayer dielectric. A deposition film of parylene can achieve a dielectric constant of 2.4, but its thermal resistance is at most about 200-300° C. and, therefore, processes for the production of semiconductor elements are restricted. Also, a porous SiO 2 film having a dielectric constant of 2.0 to 2.5 is reported, but it is poor in mechanical strength (resistance to CMP process) due to high porosity and has a problem that the pore size is not uniform. Further, these polymeric materials and porous SiO 2 film have an inferior thermal conductivity as compared to conventional SiO 2 interlayer dielectrics and accordingly may cause a problem of deterioration in wiring life (electromigration) due to rise in temperature of wiring. Use of copper as a wiring material requires covering the surface of copper wiring with a diffusion preventive film, since copper easily diffuses into insulating layers under application of an electric field. Therefore, in general, the lower and side surfaces of a copper wiring are covered with a conductive barrier metal, while the top surface thereof is covered with a silicon nitride insulating film. The dielectric constant of the silicon nitride film is about 7 and the resistance of the barrier metal is much higher than that of copper. Thus, the resistance value of the wiring as a whole increases to result in restriction on speeding up in operation of semiconductor devices. The same problem is also encountered when a low dielectric constant material is used as an insulating film. In case of using low dielectric constant insulating films, conventional silicon oxide which has a good thermal conductivity is used as a material of a layer provided with a hole for connecting the upper wiring with the lower wiring in order to avoid reduction in reliability. Since the use of this silicon oxide layer further increases wiring capacitance, a problem arises that the propagation delay of signal is caused by increase in wiring capacitance, thus resulting in restriction on speeding up of semiconductor devices. As a material having a low dielectric constant and a thermal resistance which would solve the problems as mentioned above, JP-A-2000-340689 and JP-A-2001-15496 propose low dielectric constant materials that have a borazine skeleton-based molecule in an inorganic or organic material molecule. However, the proposed low dielectric constant materials have the problem that since they are hydrolyzable, the water resistance is poor. It is an object of the present invention to provide a low dielectric constant material free from the problems as mentioned above, particularly a low dielectric constant material having an excellent water resistance as well as a low dielectric constant and a high thermal resistance. A further object of the present invention is to provide a low dielectric constant insulating film having an excellent water resistance suitable for use in semiconductor devices. A still further object of the present invention is to provide a process for preparing a low dielectric constant material having an excellent water resistance as well as a low dielectric constant and a high thermal resistance. Another object of the present invention is to provide a semiconductor device capable of operating in high speed and having a high reliability. These and other objects of the present invention will become apparent from the description hereinafter. SUMMARY OF THE INVENTION In accordance with the first aspect of the present invention, there is provided a process for preparing a low dielectric constant material comprising the step of heat-treating an inorganic or organic compound containing in its molecule a borazine skeleton structure of the formula (1-1): wherein at least one of R 1 to R 6 is a bond which binds said borazine skeleton structure to the molecule of said inorganic or organic compound, and R 1 to R 6 other than said bond are independently a hydrogen atom, an alkyl group having 1 to 20 carbon atoms, an aryl group, a substituted aryl group, an alkenyl group, an amino group, an alkylamino group, an alkoxyl group, a thioalkoxyl group, a carbonyl group, a silyl group, an alkylsilyl group, a phosphino group, an alkylphosphino group or a group of the formula: Si(OR 7 )(OR 8 )(OR 9 ) in which R 7 to R 9 are independently a hydrogen atom, an alkyl group having 1 to 20 carbon atoms, an aryl group, a substituted aryl group, an alkenyl group, an amino group, an alkylamino group, an alkoxyl group, a thioalkoxyl group, a carbonyl group, a silyl group, an alkylsilyl group, a phosphino group or an alkylphosphino group, provided that at least one of R 1 to R 6 other than said bond is not a hydrogen atom. In accordance with the second aspect of the present invention, there is provided a process for preparing a low dielectric constant material comprising the step of heat-treating a borazine skeleton-containing compound of the formula (1-2): wherein R 1 to R 6 are independently a hydrogen atom, an alkyl group having 1 to 20 carbon atoms, an aryl group, a substituted aryl group, an alkenyl group, an amino group, an alkylamino group, an alkoxyl group, a thioalkoxyl group, a carbonyl group, a silyl group, an alkylsilyl group, a phosphino group, an alkylphosphino group, or a group of the formula: Si(OR 7 )(OR 8 )(OR 9 ) in which R 7 to R 9 are independently a hydrogen atom, an alkyl group having 1 to 20 carbon atoms, an aryl group, a substituted aryl group, an alkenyl group, an amino group, an alkylamino group, an alkoxyl group, a thioalkoxyl group, a carbonyl group, a silyl group, an alkylsilyl group, a phosphino group or an alkylphosphino group, and at least one of R 1 to R 6 is not a hydrogen atom. Low dielectric constant materials having an excellent water resistance as well as a low dielectric constant and a high thermal resistance can be obtained by the above first and second processes. Thus, the present invention provides a low dielectric constant material (material I) comprising a polymeric or oligomeric, inorganic or organic material having in its molecule a borazine skeleton structure represented by any of the formulas (2) to (4): wherein R 1 to R 4 are independently a hydrogen atom, an alkyl group having 1 to 20 carbon atoms, an aryl group, a substituted aryl group, an alkenyl group, an amino group, an alkylamino group, an alkoxyl group, a thioalkoxyl group, a carbonyl group, a silyl group, an alkylsilyl group, a phosphino group, an alkylphosphino group, or a group of the formula: Si(OR 7 )(OR 8 )(OR 9 ) in which R 7 to R 9 are independently a hydrogen atom, an alkyl group having 1 to 20 carbon atoms, an aryl group, a substituted aryl group, an alkenyl group, an amino group, an alkylamino group, an alkoxyl group, a thioalkoxyl group, a carbonyl group, a silyl group, an alkylsilyl groups a phosphino group or an alkylphosphino group, provided that at least one of R 1 to R 4 is not a hydrogen atom. The present invention also provides a low dielectric constant material (material II) obtained by condensation of the compound (1-2), that is, a low dielectric constant material having a borazine skeleton-based structure formed by bonding a first borazine skeleton structure represented by any one of the formulas (2) to (4) described below with a second borazine skeleton structure represented by any one of the formulas (2) to (4) with elimination of hydrogen atoms from each of the first and second borazine skeleton structure to form a third borazine skeleton structure: wherein R 1 to R 4 are independently a hydrogen atom, an alkyl group having 1 to 20 carbon atoms, an aryl group, a substituted aryl group, an alkenyl group, an amino group, an alkylamino group, an alkoxyl group, a thioalkoxyl group, a carbonyl group, a silyl group, an alkylsilyl group, a phosphino group, an alkylphosphino group, or a group of the formula: Si(OR 7 )(OR 8 )(OR 9 ) in which R 7 to R 9 are independently a hydrogen atom, an alkyl group having 1 to 20 carbon atoms, an aryl group, a substituted aryl group, an alkenyl group, an amino group, an alkylamino group, an alkoxyl group, a thioalkoxyl group, a carbonyl group, a silyl group, an alkylsilyl group, a phosphino group or an alkylphosphino group, provided that at least one of R 1 to R 4 is not a hydrogen atom. The low dielectric constant materials of the present invention are useful as an insulating film for use in semiconductor devices, and semiconductor devices having excellent properties are obtained by the use thereof. Thus, in accordance with another aspect of the present invention, there is provided an insulating film comprising the above-mentioned low dielectric constant material I or II. In accordance with still another aspect of the present invention, there is provided a semiconductor device including such insulating film. The semiconductor devices of the present invention are operable at high speed. In a first embodiment of the semiconductor devices according to the present invention, the semiconductor device comprises a semiconductor substrate, a first insulating layer having a first trench and being formed on the surface of said semiconductor substrate, a first copper conductive layer formed so as to fill said first trench, a second insulating layer having a hole communicating with said first copper conductive layer, a second copper conductive layer filled in said hole so as to contact with said first copper conductive layer, a third insulating layer formed on said second insulating layer and having a second trench communicating with said second copper conductive layer, and a third copper conductive layer contacting said second copper conductive layer and being formed so as to fill said second trench formed in said third insulating layer, wherein at least one of said first, second and third insulating layers is made of an insulating material comprising the above-mentioned low dielectric constant material I or II. In a second embodiment of the semiconductor devices according to the present invention, the semiconductor device comprises a semiconductor substrate, a first insulating layer having a first trench and formed on the surface of said semiconductor substrate, a first copper conductive layer formed so as to fill said first trench, an insulating film which has a first hole communicating with said first copper conductive layer and which covers said first copper conductive layer and said first insulating layer, a second insulating layer having a second hole communicating with said first hole and having a second trench communicating with said second hole, a second copper conductive layer filling said first and second holes so as to contact with said first copper conductive layer, and a third copper conductive layer contacting with said second copper conductive layer and being formed so as to fill said second trench formed in said second insulating layer, wherein said insulating film is made of an insulating material comprising the above-mentioned low dielectric constant material I or II. In the above embodiments, from the viewpoints that the production of semiconductor devices is easy and the reliability of the devices is high, it is preferable that the low dielectric constant material I or II in the insulating material is amorphous. Also, from the viewpoints of excellent mechanical strength and high heat resistance, it is preferable that the low dielectric constant material I or II in the insulating material is a mixture of a microcrystalline material and an amorphous material. In the semiconductor device according to the first embodiment, from the viewpoint that the thermal conductivity is excellent and accordingly the reliability is improved, it is preferable that at least one of the first, second and third insulating layers is made of silicone oxide. Also, from the viewpoint that wiring having a good shape is obtained and accordingly the reliability is improved, it is preferable that at least one of the first, second and third insulating layers is made of an aryl ether polymer. BRIEF DESCRIPTION OF DRAWINGS FIG. 1 is a sectional view showing a wiring structure of a semiconductor device according to an embodiment of the present invention; FIG. 2 is a sectional view showing a wiring structure of a semiconductor device according to an embodiment of the present invention; FIG. 3 is a sectional view showing a wiring structure of a semiconductor device according to an embodiment of the present invention; FIG. 4 is a sectional view showing a wiring structure of a semiconductor device according to an embodiment of the present invention; FIG. 5 is a sectional view showing a wiring structure of a conventional semiconductor device; and FIG. 6 is a sectional view showing a wiring structure of a conventional semiconductor device. DETAILED DESCRIPTION The low dielectric constant materials of the present invention are prepared by subjecting a borazine derivative as a starting material, i.e., an inorganic or organic compound containing a borazine skeleton structure of the formula (1-1) in its molecule or a substituted borazine (1-2), to a condensation reaction to produce an oligomer or polymer containing the borazine skeleton structure. The preparation of the low dielectric constant materials is carried out, for example, according to the procedures as described in Yoshiharu Kimura, Senni-to-Kogyo (Fiber and Industry), Vol. 52, No. 8, 341-346 (1996); Paine & Sneddon, Recent Developments in Borazine-Based Polymers, “Inorganic and Organometallic Polymers”, American Chemical Society, 358-374 (1994); and Fazen et al., Chem. Mater., Vol. 7, p 1942 (1995). That is, the low dielectric constant materials can be obtained by heating a borazine derivative as the starting material to undergo a condensation reaction, or by firstly synthesizing a prepolymer in such a manner and then polymerizing it. In general, the condensation reaction is carried out by heating the starting material in an organic solvent at a temperature of 50 to 400° C., preferably 70 to 180° C. for 1 to 240 hours, preferably in an inert gas atmosphere such as argon. In the preparation of low dielectric constant materials is used an organic solvent which can homogeniously disperse or dissolve borazine, borazine derivatives as mentioned above or borazine-based prepolymers, e.g., an alcohol such as methanol, ethanol, propanol or butanol, acetone, benzene, toluene, xylene, glyme and others. An example of the substituted borazine (1-2) is B-triethylaminoborazine. B-triethylaminoborazine can be prepared, for example, by reacting B-trichloroborazine with ethylamine in toluene at an elevated temperature, e.g., 70° C., for several hours, e.g., 4 hours, and removing ethylamine hydrochloride and the solvent. In the inorganic or organic compound containing a borazine skeleton structure of the formula (1-1) in its molecule, the inorganic compound to which the substituted borazine (1-2) is bound includes, for instance, silicate, silazane, silsequioxane, siloxane, silane and the like. The organic compound to which the substituted borazine (1-2) is bound includes, for instance, poly(aryl ether), parylene, polyphenylene, polyphenylenevinylene, polybenzocyclobutene, polyimide, polyester, polystyrene, polymethylstyrene, polymethyl acrylate, polymethyl methacrylate, polycarbonate, adamantane, norbornene, and the like. The low dielectric constant materials of the present invention can also be obtained by a chemical vapor deposition method, as described after, using a boron source, a nitrogen source and a carbon or the like source such as methane, a chemical vapor deposition method using a substituted borazine such as methylborazine or ethylborazine, or by methods as disclosed in C. K. Narula et al., J. Am. Chem. Soc., Vol. 109, p 5556 (1987) and Y. Kimura et al., Composites Science and Technology, Vol. 51, p 173 (1994). The low dielectric constant materials of the present invention prepared from the inorganic or organic compound containing in its molecule the borazine skeleton structure shown by the formula (1-1) are inorganic or organic oligomers or polymers containing a borazine skeleton structure shown by the formula (2), (3) or (4) in the molecule thereof. These oligomers and polymers have a lower dielectric constant than silicon oxide and fluorine-containing silicon oxide, and an excellent water resistance. They are composed of, as a main component, boron nitride which has a copper diffusion preventing function and accordingly can prevent diffusion of copper. Examples of the borazine skeleton structures included in the oligomers or polymers are those having the formulas (5) to (116) shown below. The low dielectric constant materials according to another embodiment of the present invention are condensates of the substituted borazine (1-2), in other words, compounds having a third borazine skeleton-based structure formed by bonding a first borazine skeleton structure represented by any one of the formulas (2) to (4) with a second borazine skeleton structure represented by any one of the formulas (2) to (4) with elimination of hydrogen atoms from each of the molecules of a substituted borazine to form the third borazine skeleton structure. Examples of the condensates are, for instance, compounds having borazine skeletons structures shown by the above formulas (25) to (28). The reason why the low dielectric constant material of the present invention can achieve a low dielectric constant is considered that the electronic polarization is decreased by an ionic electronic structure of the borazine skeleton. Also, a high heat resistance can be achieved by the low dielectric constant materials of the present invention, since inorganic polymeric materials which have of course a higher heat resistance than organic polymeric materials are used. Further, the reason why the low dielectric constant materials of the present invention have a high water resistance is considered that if R 1 to R 4 is substituents other than hydrogen atom in the formulas (2) to (4), they firmly bond to boron atom or nitrogen atom in the borazine skeleton and are prevented from reacting with water. Since a hydrogen atom bonding to a boron atom or a nitrogen atom is easily hydrolyzed, it is necessary that in the low dielectric constant material of the present invention, at least one of R 1 to R 4 in the formulas (2) to (4) is not a hydrogen atom, but a substituent. In particular, since a hydrogen atom bonding to a boron atom causes a hydrolysis reaction more easily as compared with that bonding to a nitrogen atom, it is preferable that a substituent is bonded to a boron atom. As to the degree of substitution, preferred from the viewpoint of water resistance, of hydrogen atoms on the borazine skeletons included in a molecule which constitutes the low dielectric constant material, assuming that the degree of substitution is 100% if all hydrogen atoms on the borazine skeletons are substituted by a substitutent or substituents shown in the formulas (2) to (4), water resistance equivalent to that for a degree of substitution of 100% is obtained when 30 to 40% of all hydrogen atoms are substituted by a substitutent or substituents shown in the formulas (2) to (4), namely when the degree of substitution is 30 to 40%. The dielectric constant can be further lowered by introducing fluorine atom (F) into boron nitride. Thus, an insulation layer having a lower dielectric constant can be obtained thereby. The insulating films of the present invention are obtained by forming the low dielectric constant materials of the present invention into thin films. The insulating films of the present invention are applicable as an interlayer insulating film of semiconductor devices, whereby excellent semiconductor devices can be obtained. In case of using the low dielectric constant materials in the form of a film, for example, as an interlayer insulating film for semiconductor devices, the film can be formed by coating a solution or dispersion of the low dielectric constant material in a solvent. In that case, the low dielectric constant material may be used in combination with other materials such as other insulating materials which are used preferably in an amount of at most 20% by weight based on the total weight of the low dielectric constant material of the present invention and other materials. Examples of the other materials are, for instance, a known interlayer insulating material for semiconductor devices such as silicate, silazane, silsequioxane, siloxane, silane, polyaryl ether, parylene or polybenzocyclobutadiene, a general insulating material such as adamantane, norbornene, polyimide, polyester, polystyrene, polymethylstyrene, polymethyl acrylate, polymethyl methacrylate or polycarbonate, an amine such as cyclohexylamine, aniline or ethylamine, a surface active agent, and the like. The coating to a substrate can be conducted by spray coating, dip coating, spin coating or other known coating methods. The solvent or dispersing medium includes, for instance, acetone, benzene, glyme, tetrahydrofuran, chloroform and other organic solvents capable of dissolving or dispersing the low dielectric constant materials. The concentration is preferably from 10 to 30% by weight. Preferably, after drying the coated film, the dried film is further heat-treated to cure the film at a temperature of 300 to 450° C., preferably 350 to 400° C. The thickness of the insulating film is preferably from 0.3 to 0.8 μm. In case of using the low dielectric constant materials as a film such as an interlayer insulating film for semiconductor devices, thin films can also be formed according to procedures as described for example in S. V. Nguyen, T. Nguyen, H. Treichel and O, Spindler, J. Electrochem. Soc., Vol. 141, No. 6, 1633-1638 (1994); W. F. Kane, S. A. Cohen, J. P. Hummel and B. Luther, J. Electrochem. Soc., Vol. 144, No. 2, 658-663 (1997); and M. Maeda and T. Makino, Japanese Journal of Applied Physics, Vol. 26, No. 5, 660-665 (1987). For example, the insulating film or layer can be obtained by subjecting a mixture of diborane (B 2 H 6 ), ammonia (NH 3 ) and methane or a mixture of borazine (B 3 H 3 N 6 ), nitrogen (N 2 ) and methane as a raw material a chemical vapor deposition method (CVD method), thereby causing a condensation reaction. In case that the low dielectric constant materials are used in the form of a bulk body as a low dielectric constant substrate, the materials are molded by casting into a mold and heat-treating the resulting molded article. The low dielectric constant material to be cast may be used in combination with other materials as mentioned above. The content of other materials is at most 20% by weight. The insulating films of the present invention applicable to various electronic parts as an interlayer insulating film for semiconductor devices, as a barrier metal layer or etch stopper layer, and as an IC substrate. Thus, the present invention provides semiconductor devices including an insulating layer or film made of the low dielectric constant materials of the present invention. In an embodiment of the semiconductor devices according to the present invention, a first insulating layer having a first copper conductive layer disposed to form a lower wiring and a third insulating layer having a third copper conductive layer disposed to form an upper wiring are stacked on the surface of a semiconductor substrate through a second insulating layer interposed therebetween and having a second copper conductive layer communicating with both the first copper conductive layer and the third copper conductive layer so as to electrically connect the lower wiring with the upper wiring. In this embodiment, at least one of the first, second and third insulating layers is made of an insulating material containing the low dielectric constant material of the present invention. In another embodiment of the semiconductor devices according to the present invention, a first insulating layer having a first copper conductive layer disposed to form a lower wiring and a second insulating layer having a third copper conductive layer disposed to form an upper wiring and having a second copper conductive layer communicating with both the first copper conductive layer and the third copper conductive layer so as to electrically connect the lower wiring with the upper wiring are stacked on the surface of a semiconductor substrate through an insulating film interposed therebetween, the second copper conductive layer also extending through the insulating film. In this embodiment, the insulating film interposed between the first and second insulating layers is made of an insulating material containing the low dielectric constant material of the present invention. Since the insulating layer or film made of an insulating material containing the low dielectric constant material of the present invention is used in the above semiconductor devices instead of conventional built-up films of silicon oxide and silicon nitride, the wiring capacitance can be reduced. Also, since the insulating layer or film is made of an insulating material containing the low dielectric constant material of the present invention which has a copper diffusion preventing function, it is not needed to use a barrier metal layer at connecting hole portions and, therefore, a low resistant wiring can be obtained and it is possible to operate the semiconductor devices at high speed. In the above embodiments, the first, second and third conductive layers are made of copper and, therefore, the wiring delay can be decreased as compared with the use of aluminum, but the materials of the conductive layers are not limited copper. An example of the wiring structure of semiconductor devices according to the present invention is shown in FIG. 1 . In the figure, numeral 1 denotes a semiconductor substrate made of silicon, and numeral 19 denotes an insulating layer made of silicon oxide. On the silicon oxide insulating layer 19 is formed an insulating layer 20 having a thickness of 0.3 μm and made of a crosslinked poly(B-methylaminoborazine) which is a low dielectric constant material according to the present invention. The insulating layers 19 and 20 constitute the first insulating layer. In the insulating layer 20 is formed a first trench 3 having a width of 0.2 μm and a depth of 0.2 μm in the pattern of a first wiring. A first copper conductive layer 5 is filled in the trench 3 . A second insulating layer 21 having a thickness of 0.5 μm made of the crosslinked poly(B-methylaminoborazine) is formed on the insulating layer 20 and the first copper conductive layer 5 . In the second insulating layer 21 is formed a hole 8 having a diameter of 0.15 μm and extending to the first copper conductive layer 5 , and the hole 8 is filled with copper to form a second copper conductive layer 10 so as to contact the first copper conductive layer 5 . On the insulating layer 21 is formed a third insulating layer 22 having a thickness of 0.2 μm made of the crosslinked poly(B-methylaminoborazine). In the third insulating layer 22 is formed a second trench 12 having a depth of 0.2 μm in the pattern of a second wiring. The bottom of the trench 12 extends to the insulating layer 21 , and copper is filled in the trench 12 to form a third copper conductive layer 13 . An insulating film 23 made of the crosslinked poly(B-methylaminoborazine) is formed on the insulating layer 22 and the third copper conductive layer 13 . In semiconductor devices having such a structure, all copper conductive layers, that is, the first copper conductive layer 5 , the second copper conductive layer 10 and the third copper conductive layer 13 , are in contact with the insulating layers 20 , 21 and 22 and film 23 made of an insulating material comprising the low dielectric constant material of the present invention. Thus, copper diffusion from the conductive layers can be prevented from occurring. Furthermore, since the insulating layers 20 , 21 , 22 and 23 have a dielectric constant of 2.2 and also do not require a barrier metal layer, the wiring capacitance can be reduced as compared with conventional wiring structure shown in FIG. 6 , whereby high speed operation of semiconductor devices can be ensured. FIG. 2 is a sectional view of a semiconductor device showing a further embodiment of the present invention. An insulating layer 19 made of silicon oxide is formed on a silicon semiconductor substrate 1 . On the silicon oxide insulating layer 19 is formed an insulating layer 20 a having a thickness of 0.3 μm and made of an amorphous crosslinked poly(B-methylaminoborazine) which is a low dielectric constant material according to the present invention. The insulating layers 19 and 20 a constitute the first insulating layer. In the insulating layer 20 a is formed a first trench 3 having a width of 0.2 μm and a depth of 0.2 μm in the pattern of a first wiring. A first copper conductive layer 5 is filled in the trench 3 . A second insulating layer 21 b having a thickness of 0.5 μm made of a mixture of microcrystalline and amorphous crosslinked poly(B-methylaminoborazine) is formed on the insulating layer 20 a and the first copper conductive layer 5 . In the second insulating layer 21 b is formed a hole 8 having a diameter of 0.15 μm and extending to the first copper conductive layer 5 , and the hole 8 is filled with copper to form a second copper conductive layer 10 so as to contact the first copper conductive layer 5 . On the insulating layer 21 b is formed a third insulating layer 22 a having a thickness of 0.2 μm made of the same material as the insulating layer 20 a , namely amorphous crosslinked poly(B-methylaminoborazine). In the third insulating layer 22 a is formed a second trench 12 having a depth of 0.2 μm in the pattern of a second wiring. The bottom of the trench 12 extends to the insulating layer 21 b , and copper is filled in the trench 12 to form a third copper conductive layer 13 . An insulating film 23 b made of the same material as the insulating layer 21 b is formed on the insulating layer 22 a and the third copper conductive layer 13 . In semiconductor devices having such a structure, all copper conductive layers, that is, the first copper conductive layer 5 , the second copper conductive layer 10 and the third copper conductive layer 13 , are in contact with the insulating layers 20 , 21 and 22 and film 23 made of an insulating material comprising the low dielectric constant material of the present invention. Thus, copper diffusion from the conductive layers can be prevented from occurring. Furthermore, since the insulating layers 20 , 21 , 22 and 23 have a dielectric constant of 2.3 and also do not require a barrier metal layer, the wiring capacitance can be reduced as compared with conventional wiring structure shown in FIG. 6 , whereby high speed operation of semiconductor devices can be ensured. FIG. 3 is a sectional view of a semiconductor device showing another embodiment of the present invention. An insulating layer 19 made of silicon oxide is formed on a silicon semiconductor substrate 1 . On the silicon oxide insulating layer 19 is formed an insulating layer 25 having a thickness of 0.2 μm and made of a poly(aryl ether). The insulating layers 19 and 25 constitute the first insulating layer. In the insulating layer 25 is formed a first trench 3 having a width of 0.2 μm and a depth of 0.2 μm in the pattern of a first wiring. A first copper conductive layer 5 is filled in the trench 3 . A first conductive film (barrier metal film) 4 having a diffusion preventive function is formed so as to cover the surface of the trench 3 . The barrier metal film 4 is made of tantalum nitride and has a thickness within the range of 10 to 20 nm. Copper is filled in the trench 3 covered with the barrier metal film 4 to form a first copper conductive layer 5 . A second insulating layer 21 b having a thickness of 0.5 μm made of a mixture of microcrystalline and amorphous crosslinked poly(B-methylaminoborazine), which is the low dielectric constant material of the present invention, is formed on the insulating layer 25 and the first copper conductive layer 5 . In the second insulating layer 21 b is formed a hole 8 having a diameter of 0.15 μm and extending to the first copper conductive layer 5 , and the hole 8 is filled with copper to form a second copper conductive layer 10 so as to contact the first copper conductive layer 5 . On the insulating layer 21 b is formed a third insulating layer 27 made of the same material as that of the insulating layer 25 , i.e., poly(aryl ether), and having a thickness of 0.2 μm. In the third insulating layer 27 is formed a second trench 12 having a depth of 0.2 μm in the pattern of a second wiring. The bottom of the trench 12 extends to the insulating layer 21 b . A second conductive film (barrier metal film) 11 having a diffusion preventive function against copper is formed so as to cover the inner surface of the trench 12 . The barrier metal film 11 has the same composition and the same thickness as those of the barrier metal film 4 . Copper is filled in the trench 12 covered with the barrier metal film 11 to form a third copper conductive layer 13 . An insulating film 23 b made of the same material as the insulating layer 21 b is formed on the insulating layer 27 and the third copper conductive layer 13 . In semiconductor devices having such a structure, the first copper conductive layer 5 is in contact with the barrier metal film 4 and the insulating layer 21 b , and the third copper layer 13 is in contact with the barrier metal film 11 and the insulating layer 23 b . Further, the second copper conductive layer 10 is in contact with the barrier metal 11 and the insulating layer 21 b . Because of having such a structure, diffusion of copper from the conductive layers can be prevented. Moreover, since the insulating layers 25 and 27 made of poly(aryl ether) have a dielectric constant of 2.8 and the insulating layers 21 b and 23 b made of crosslinked poly(B-methylaminoborazine) have a dielectric constant of 2.2, the wiring capacitance can be reduced to a level lower than that achieved by a conventional wiring structure shown in FIG. 6 , whereby a high speed operation of semiconductor devices is made possible. Further, since the insulating layers 25 and 27 are made of poly(aryl ether) and the insulating layers 21 b and 23 b are made of crosslinked poly(B-methylaminoborazine), the etching selective ratio is high and accordingly it is possible to form wiring having a good shape. In this embodiment, the layer in which second copper conductive layer 10 is provided, i.e., insulating layer 21 b , is formed from a crosslinked poly(B-methylaminoborazine). Substantially the same effect can be obtained also when the layer provided with the first or third copper conductive layer 5 or 13 , i.e., insulating layer 25 or 27 , is formed from the crosslinked poly(B-methylaminoborazine). Another example of the wiring structure of semiconductor devices using the low dielectric constant material of the present invention as an insulating film or layer is shown in FIG. 4 . A first insulating layer 29 made of silicon oxide is formed on a silicon semiconductor substrate 1 . In the insulating layer 29 is formed a trench 3 having a width of 0.2 μm and a depth of 0.2 μm in the pattern of a first wiring. A first conductive film (barrier metal film) 4 having a diffusion preventive function is formed so as to cover the surface of the trench 3 . The barrier metal film 4 is made of tantalum nitride and has a thickness within the range of 10 to 20 nm. Copper is filled in the trench 3 covered with the barrier metal film 4 to form a first copper conductive layer 5 . An insulating layer 30 b having a thickness of 0.05 μm made of a mixture of microcrystalline and amorphous crosslinked poly(B-methylaminoborazine), in other words, microcrystals-containing amorphous crosslinked poly(B-methylaminoborazine), which is the low dielectric constant material of the present invention, is formed on the insulating layer 29 and the first copper conductive layer 5 . On the insulating layer 30 b is formed a second insulating layer 31 made of silicon oxide. In the second insulating layer 31 are formed a hole 8 having a diameter of 0.15 μm and a trench 12 having a depth of 0.2 μm and a second wiring pattern. The hole 8 extends from the first conductive layer 5 to the trench 12 formed in the surface region of the insulating layer 31 through the insulating layer 30 b and the insulating layer 31 . Second and third conductive films (barrier metal films) 9 and 11 made of tantalum nitride having a diffusion preventive function are formed so as to cover the surfaces of the hole 8 and the trench 12 . Copper is filled in the hole 8 and the trench 12 to form second copper conductive layer 10 and third copper conductive layer 13 , respectively. The barrier metal film is also formed at the interface between the first copper conductive layer 5 and the second copper conductive layer. An insulating film 23 b made of the same material as the insulating layer 30 b is formed on the insulating layer 31 and the third copper conductive layer 13 . In semiconductor devices having such a structure, the first, second and third copper conductive layers 5 , 10 and 13 are in contact with the barrier metal films 4 , 9 and 11 and the insulating layers 23 b and 30 b . Thus, diffusion of copper from the conductive layers 5 , 10 and 13 can be prevented. Moreover, since the insulating layers 23 b and 30 b have a dielectric constant of 2.2 and the insulating layers 29 and 31 have a dielectric constant of 4.2, the wiring capacitance can be reduced to a level lower than that achieved by a conventional wiring structure shown in FIG. 6 , whereby a high speed operation of semiconductor devices is made possible. The present invention is more specifically described and explained by means of the following examples. EXAMPLE 1 Soluble poly(B-trimethylborazilene) was synthesized according to Fazen et al's method disclosed in Fazen et al., Chem. Mater., Vol. 7, p 1942, 1995. Tetraglyme was used as a solvent, and B-trimethylborazine was heated in an Ar gas at 220° C. for two weeks with stirring and degassing to give a highly viscous liquid. The liquid was evaporated to give a white powder of a low dielectric constant material according to the present invention. This material had a chemical structure shown by the following formula (117): and had an average molecular weight (Mn) of about 7,500. The obtained low dielectric constant material was dissolved in acetone and coated by spin coating method onto a quartz plate on which gold was deposited to form a counter electrode. The coated film was then dried at 100° C. for 10 minutes and heated at 400° C. for 10 minutes to give an insulating film according to the present invention. The thus heat-treated film was made of a partially crosslinked poly(B-methylboradine). Gold was deposited onto the obtained insulating film as a main electrode. EXAMPLE 2 Synthesis of soluble poly(B-triethylborzilene) was carried out in the same manner as Example 1. Tetraglyme was used as a solvent, and B-triethylborazine was heated in an Ar gas at 220° C. for two weeks with stirring and degassing to give a highly viscous liquid. The liquid was evaporated to give a white powder of a low dielectric constant material according to the present invention. This material had a chemical structure shown by the following formula (118): and had an average molecular weight (Mn) of about 5,500. An insulating film was formed from the obtained low dielectric constant material by conducting the spin coating in the same manner as in Example 1 and drying at 100° C. for 10 minutes. Gold was then deposited onto the insulating film as a main electrode. EXAMPLE 3 A white powder of poly(methylborazinylamine) was prepared according to Narula et al's method disclosed in C. K. Narula, R. Schaeffer, R. T. Paine, A. K. Datye and W. F. Hammetter, J. Am. Chem. Soc., Vol. 109, p 5556 (1987). The thus obtained low dielectric constant material was dispersed into acetone, and the dispersion was coated by spin coating and dried at 100° C. for 10 minutes in the same manner as in Example 1 to give an insulating film. Gold was then deposited thereon as a main electrode. EXAMPLE 4 A white powder of poly(B-methylaminoborazine) was prepared according to Kimura et al's method disclosed in Y. Kimura et al., Composites Science and Technology, Vol. 51, p 173 (1994). The thus obtained low dielectric constant material was dispersed into acetone, and the dispersion was coated by spin coating and dried at 100° C. for 10 minutes in the same manner as in Example 1 to give an insulating film. Gold was then deposited thereon as a main electrode. Dielectric constants of the insulating films obtained in Examples 1 to 4 were measured at 25° C. and 1 MHz by using an impedance analyzer (model 4191A made by Hewlett Packard). In order to evaluate the water resistance, the dielectric constant was also measured with the lapse of time. The results are shown in Table 1. COMPARATIVE EXAMPLE An insulating film was formed from polyboradilene in the same manner as in Example 1, and the dielectric constant thereof was measured. The result is shown in Table 1. TABLE 1 Dielectric constant After preparation After After After of sample 1 day 2 days 3 days Example 1 2.0 2.1 2.2 2.2 Example 2 2.1 2.1 2.1 2.1 Example 3 2.2 2.2 2.2 2.2 Example 4 2.1 2.1 2.1 2.2 Com. Ex. 2.0 3.5 4.2 4.5 The insulating films obtained in Examples 1 to 4 have a dielectric constant of at most 2.4. From these results, it is understood that a substrate having a low dielectric constant can be obtained. Also, these polymeric borazine compounds can be graphitized by heating at a temperature of 1,000 to 1,200° C. (Application View of Inorganic Polymer, p 70, 1990, supervised by Naruyuki Kajiwara). Thus, these insulating films have a thermal resistance of at least 450° C. Further, as apparent from the results shown in Table 1, the films obtained in Examples 1 to 4 show no or little change in dielectric constant with the lapse of time. Thus, it is understood that these films have an excellent water resistance.

A low dielectric constant material having an excellent water resistance obtained by heat-treating a borazine compound of the formula (1-2): or an inorganic or organic compound having a group derived from the borazine compound (1-2) to undergo a condensation reaction, thereby producing an oligomer or polymer, wherein R 1 to R 6 are independently a hydrogen atom, an alkyl group having 1 to 20 carbon atoms, an aryl group, a substituted aryl group, an alkenyl group, an amino group, an alkylamino group, an alkoxyl group, a thioalkoxyl group, a carbonyl group, a silyl group, an alkylsilyl group, a phosphino group, an alkylphosphino group, or a group of the formula: Si(OR 7 )(OR 8 )(OR 9 ), and at least one of R 1 to R 6 is not hydrogen atom.

FIELD OF THE INVENTION The present invention relates generally to digital image compression. BACKGROUND OF THE INVENTION Several techniques exist for compressing digital image files. Compression is done in order to reduce the resource requirements for storing and transmitting image files. Lossless compression techniques exploit statistical characteristics of the image data to code the files more efficiently, and allow exact reconstruction of the original data. “Lossy” compression techniques use similar statistical methods, and also tolerate small changes in the content of the files after compression and reconstruction. Lossy techniques typically produce compressed files considerably smaller than files produced by lossless techniques, and in some applications, the changes in content are negligible. One commonly used lossy compression technique is the JPEG technique, named for the Joint Photographic Experts Group, the committee that developed the specifications for standard use of the technique and for the standard file format of JPEG image files. The JPEG technique is especially useful for images of natural scenes, and is widely used for compressing digital photographs. Many digital still cameras include circuitry that implements the JPEG standard to create compressed files. Some digital cameras provide the ability to capture moving pictures as well as still images. Moving pictures may be thought of as sequences of still images. To facilitate the compression of moving pictures, another standard, called MPEG, has been developed by the Moving Picture Experts Group. There are several variants of MPEG compression, but the features described in this specification are common to all, so all the variants will be referred to here generically as MPEG. In a simple implementation of MPEG, a moving picture sequence comprises a series of individually compressed still images called “I-frames”. An MPEG I-frame is intra-coded, that is, compressed without regard to the content of frames occurring before or after it in the sequence. The MPEG technique allows, but does not require, other kinds of frames, for example “P-frames” and “B-frames”, that do take into account the content of adjoining frames. The present invention addresses the generation of I-frames. Some of the processing necessary to construct an MPEG I-frame is identical to some of the processing used to construct a JPEG compressed image. However, a finishing step is significantly different between the two techniques. The circuitry or other engine used in cameras to construct JPEG images is often configurable in order to allow the compression to be optimized for particular data, and some flexibility is allowed within the JPEG specification. However, it is not possible to construct a completed MPEG I-frame using a standard JPEG engine or circuitry. A brief and simplified example will aid in providing an overview of the steps involved in JPEG and MPEG compression. An MPEG I-frame has many similarities to a still image compressed using the JPEG technique. The sequence of steps required for generating either a JPEG image or an MPEG I-frame includes: 0. Color space conversion 1. Downsampling, also called subsampling or decimation 2. Constructing macroblocks 3. Performing a Discrete Cosine Transform (DCT) 4. Quantization 5. “Zig zag” ordering of the quantized coefficients 6. Differential coding of the DC coefficient from the DCT 7. Run-length coding of the AC coefficients from the DCT 8. Variable-length coding of the coefficients from the DCT All of these steps except the last may be performed identically whether the desired result is a JPEG image or an MPEG I-frame. However, the final step of variable-length coding the coefficients is significantly different for constructing an MPEG I-frame than for constructing a JPEG image. A digital camera produces an ordered array of data representing an original scene. Each location in the scene is represented by a corresponding picture element, or “pixel”. The data describing each pixel indicate the brightness and color of the original scene at the location corresponding to the pixel. The brightness and color are often represented by numerical values indicating the strengths of red, green, and blue light sensed from the scene location. An image of this type is often said to be in “RGB” format. Other representations of brightness and color may be used, and conversions from one system of representation, or “color space”, to another are readily accomplished. Both JPEG and MPEG require the image to be represented in the color space known as YCrCb. In the YCrCb color space, a pixel is described by its overall brightness or luminance, (Y) and two chrominance values (Cr and Cb) that describe the color of the pixel. The color space conversion step of JPEG or MPEG compression involves converting from another color space such as RGB to YCrCb. Many cameras use electronic array sensors that have many more pixels than are typically used in moving picture frames. Often, cameras provide the ability to save images at various resolutions. The lower the resolution, the fewer pixels used to represent the image and the less detail will be visible in the image file. The conversion from a higher resolution image to a lower resolution image is often called downsampling, subsampling, or decimation. Additionally, the MPEG specification requires and the JPEG specification allows the chrominance channels of the image to be further downsampled in the 4:2:0 video format. In this format, the chrominance channels are downsampled to half the linear resolution of the luminance channel in each of the two orthogonal coordinate directions of the image. Thus each chrominance channel represents the image with one fourth as many pixels as does the luminance channel, and at a correspondingly lower resolution. Chrominance downsampling takes advantage of the human visual system's decreased sensitivity to resolution in the chrominance channels in comparison with the luminance channel to reduce the data required to represent a pleasing image. Once the image is downsampled, it is divided into “macroblocks”. A macroblock consists of a 16-pixel by 16-pixel sample array of luminance samples together with one 8-sample by 8-sample block from each of the chrominance channels. The sample array of luminance samples may be thought of as four subarrays that are each eight pixels square. Images that are not a multiple of 16 pixels wide or tall are padded with blank pixels so that complete macroblocks may be constructed. The next step in the process uses the data in arrays of numbers eight elements square. The division of the image into macroblocks may be entirely conceptual, as the data in the memory of the camera, imaging device, or system need not be rearranged to accomplish the division. Identifying the macroblocks partitions all of the image data, both luminance and chrominance, into arrays that are eight elements square. For example, an array of luminance samples may be as follows: 102 100 101 101 104 104 122 137 (1) 102 100 100 101 104 108 121 132 104 102 101 101 105 106 123 135 107 105 103 99 107 109 123 134 110 105 104 104 109 110 126 138 112 109 107 97 111 113 129 139 114 102 113 112 122 121 136 153 124 118 124 124 140 151 164 181 (1) This example array of luminance data will be used below to illustrate the following steps, and to describe an embodiment of the invention. One of ordinary skill in the art will recognize that the steps and the embodiment of the invention apply to both luminance and chrominance data, and that no loss of generality is intended or created by using a single example array. For each 8×8 array in the image, a two-dimensional discrete cosine transform (DCT) is performed. The DCT is described in MPEG Video Compression Standard , edited by Joan L. Mitchell, William B. Pennebaker, Chad E. Fogg, and Didier J. LeGall, and published by Chapman & Hall, ISBN 0-412-08771-5. The DCT of the example array above is: 928.12 −86.29 53.66 −15.12 13.12 −3.35 1.18 11.27 (2) −64.23 18.27 −2.00 −5.23 −1.06 1.39 −5.46 −4.22 36.50 −18.85 −1.66 −1.36 2.67 .89 3.53 −.37 −25.06 11.06 1.78 −1.51 .19 −.14 −1.19 2.27 19.38 −6.59 1.41 .14 −.13 −.72 −.18 −1.64 −11.01 3.31 −.84 −2.72 2.88 .39 .76 2.63 6.12 −1.25 4.78 .60 −3.68 −2.55 −1.84 .77 −1.07 −1.29 −1.92 −3.46 5.36 3.18 −.24 −.65 The upper left DCT coefficient indicates a scaled average value of the input data array. In general, the other coefficients represent the spatial frequency content of the image, with higher frequency components at the lower right of the array. The next step in both JPEG and MPEG compression is to “quantize” the array. Quantization is performed by an element-by-element division by another array of quantizing values, and rounding the results. An example array of quantizing values is: 8 16 19 22 26 27 29 34 (3) 16 16 22 24 27 29 34 37 19 22 26 27 29 34 34 38 22 22 26 27 29 34 37 40 22 26 27 29 32 35 40 48 26 27 29 32 35 40 48 58 26 27 29 34 38 46 56 69 27 29 35 38 46 56 69 83 Using array (3) to quantize the array (2) of DCT coefficients above gives these quantized coefficients: 116 −5 2 0 0 0 0 0 (4) −4 1 0 0 0 0 0 0 1 0 0 0 0 0 0 0 −1 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 0 After the quantization, the coefficients are placed in a “zig zag” order. The order of reading out the coefficients is illustrated below: 1 2 6 7 15 16 28 29 (5) 3 5 8 14 17 27 30 43 4 9 13 18 26 31 42 44 10 12 19 25 32 41 45 54 11 20 24 33 40 46 53 55 21 23 34 39 47 52 56 61 22 35 38 48 51 57 60 62 36 37 49 50 58 59 63 64 Because coefficients in the lower right part of the array are likely to be zero after quantization, the zig zag ordering tends to maximize runs of zeros in the ordered list. The coefficients of the example array in zig zag order are: 116 −5 −4 1 1 2 0 0 0 −1 0 0 0 0 . . . (50 more zeros) The first coefficient in this list represents a scaled average value for the pixels in the 8×8 block. This is often called the “DC” coefficient. The other coefficients are called “AC” coefficients. In both JPEG and MPEG, the DC coefficient for each block is differentially coded. That is, rather than store the coefficient itself, the difference between the coefficient from the previous block and the present coefficient is stored. Because the DC coefficients tend to change slowly, this differential coding tends to allow the storage of smaller numbers, thereby conserving storage space. In this example, it is assumed that the previous pixel block had a DC coefficient after quantization of 120, resulting in a difference of −4. The coefficients can be further arranged as follows: TABLE 1 Coefficient Preceding Number run of zeros Value 0 (DC) N/A −4 1 0 −5 2 0 −4 3 0 1 4 0 1 5 0 2 9 3 −1 End of Block The final step in compressing a block of pixels is to encode this information using variable length coding, which is often called Huffman coding. In Huffman coding, common patterns in data are assigned short sequences of bits, while less common patterns are assigned longer sequences. The sequences are chosen so that they cannot be confused with each other. In this way, data that have a nonuniform distribution of pattern frequencies can be stored losslessly in a smaller form. In both JPEG and MPEG, different codes are used for the DC and AC coefficients. In JPEG, different codes may be used for the luminance channel and the chrominance channels. MPEG specifies the table of Huffman codes for the quantized DCT coefficients. JPEG allows the user to select a set of Huffman codes. It is possible to select the JPEG codes for the DC coefficient to match the MPEG specification. However, the coding schemes are significantly different between JPEG and MPEG for the AC coefficients, and it is not possible to configure a JPEG engine to generate the Huffman code stream of an MPEG file. The Huffman codes for the DC coefficient of the luminance channel for an MPEG file and a typical JPEG file are selected according to the following table: TABLE 2 Y code size magnitude range 100 0 0 00 1 −1, 1 01 2 −3 . . . −2, 2 . . . 3 101 3 −7 . . . −4, 4 . . . 7 110 4 −15 . . . −8, 8 . . . 15 1110 5 −31 . . . −16, 16 . . . 31 11110 6 −63 . . . −32, 32 . . . 63 111110 7 −127 . . . −64, 64 . . . 127 1111110 8 −255 . . . −128, 128 . . . 255 In this table, the Y code is a bit pattern that identifies the size range of a particular DC coefficient. The “size” entry indicates the number of bits that follow the Y code to indicate the exact value of the coefficient. The magnitude range indicates the values represented by various bit patterns. In the example from above, the value to be encoded is −4. The fourth line in the table encompasses a value of −4, so the Y code bit pattern to be used is 101 . The table indicates that a three-bit value follows this Y code. There are eight possible patterns of three bits, and there are eight values in the table that correspond to the eight patterns— −7, −6, −5, −4, 4, 5, 6, and 7. The bit patterns corresponding to these values are as follows: TABLE 3 Bit pattern value 000 −7 001 −6 010 −5 011 −4 100 4 101 5 110 6 The bit pattern column in this table is simply the possible bit patterns in ascending order, and the corresponding values are the possible values in ascending order. Similar tables can be constructed for other lines in Table 2. From Table 3, the following bit pattern for a value of −4 is 011. The value stored in the file to indicate a DC coefficient of −4 is then 101 011. Thus six bits are required to represent this value. By way of further example, a coefficient value of −1 would be represented by a bit pattern of 00 0, requiring only three bits. A coefficient value of 129 would be represented by a bit pattern of 1111110 10000001, requiring 15 bits. Because the DC coefficients tend to be small, most require only a small number of bits for representation, resulting in efficient storage of the DC coefficients in an image. The JPEG specification provides for generating Table 2 algorithmically from a list of the number of codes of each size to be generated and an ordering of the categories represented by the codes. The list of the number of codes of each size is an array of 16 numbers, and the ordering is an array containing as many numbers as the total of the entries in the first list. The arrays needed to generate Table 2 are: Code length counts: 0, 2, 3, 1, 1, 1, 1, 0, 0, 0, 0, 0, 0, 0, 0, 0 Ordering values: 1, 2, 0, 3, 4, 5, 6, 7, 8 Specifying these arrays to the JPEG circuitry or other engine in a camera or system serves to “configure” the circuitry or engine. The arrays may be specified for each image, and are stored in the resulting file with the image data so that the data may be reconstructed. The JPEG technique allows the arrays to be specified for each image so that the Huffman codes may be optimized for maximum compression if the programmer so desires. A computer program for generating the code tables from the arrays is given in Appendix A. A table similar to Table 2 may be constructed for the chrominance channel DC coefficients of the image, using uses different generating values, and resulting in different Huffman codes for the values. Coding of the AC coefficients is done differently than the DC coefficients. JPEG and MPEG also code the AC coefficients differently from each other. JPEG AC coding is discussed first below. In a JPEG file, Huffman codes are assigned not just to the size range of the AC coefficient, but a combination of the coefficient size range and the number of zero coefficients preceding the non-zero coefficient. For example, a coefficient may have a value of 9 and follow a run of 3 zeros. This coefficient is said to have a run/size combination of 3/4. A coefficient with a value of −1 and following another non-zero coefficient would have a run/size combination of 0/1. Each run/size combination is assigned a Huffman code. Each non-zero AC coefficient is represented in the resulting JPEG file by its proper Huffman code (indicating the number of zero coefficients preceding the non-zero coefficient and the relative size of the non-zero coefficient) and a set of following bits that specify the exact value of the coefficient. The following bits for the AC coefficients are as described for the DC coefficient in Tables 2 and 3. A typical JPEG table for coding AC coefficients (analogous to Table 2 above for coding DC coefficients) is abbreviated below: TABLE 4 Run/size Code 0/0 1010 (Special end-of-block character) 0/1 00 0/2 01 0/3 100 0/4 1011 0/5 11010 . . . 1/1 1100 1/2 11011 1/3 1111011 . . . 2/1 11100 2/2 11111001 . . . 3/1 111010 3/2 111110111 . . . Combining tables 1, 2, 3, and 4 above, it is now possible to determine the JPEG bit pattern for the luminance values of the entire example pixel block: TABLE 5 (JPEG) Coefficient Number Run/size Value Bit pattern 0 (DC) N/A −4 101 011 1 0/3 −5 100 010 2 0/3 −4 100 011 3 0/1 1 00 1 4 0/1 1 00 1 5 0/2 2 01 10 9 3/1 −1 111010 0 End of Block 1010 The JPEG specification provides for generating Table 4 algorithmically from a list of the number of codes of each size to be generated and an ordering of the categories represented by the codes, in the same way that Table 2 can be generated. The arrays needed to generate Table 4 are: Code length counts: 0, 2, 1, 3, 3, 2, 4, 3, 5, 5, 4, 4, 0, 0, 1, 125 Ordering values (in hexadecimal notation):  01, 02, 03, 00, 04, 11, 05, 12, 21, 31, 41, 06, 13, 51, 61, 07 22, 71, 14, 32, 81, 91, A1, 08, 23, 42, B1, C1, 15, 52, D1, F0 24, 33, 62, 72, 82, 09, 0A, 16, 17, 18, 19, 1A, 25, 26, 27, 28 29, 2A, 34, 35, 36, 37, 38, 39, 3A, 42, 44, 45, 46, 47, 48, 49 4A, 53, 54, 55, 56, 57, 58, 59, 5A, 63, 64, 65, 66, 67, 68, 69 6A, 73, 74, 75, 76, 77, 78, 79, 7A, 83, 84, 85, 86, 87, 88, 89 8A, 92, 93, 94, 95, 96, 97, 98, 99, 9A, A2, A3, A4, A5, A6, A7 A8, A9, AA, B2, B3, B4, B5, B6, B7, B8, B9, BA, C2, C3, C4, C5 C6, C7, C8, C9, CA, D2, D3, D4, D5, D6, D7, D8, D9, DA, E1, E2 E3, E4, E5, E6, E7, E8, E9, EA, F1, F2, F3, F4, F5, F6, F7, F8 F9, FA In the above array of ordering values, the first hex digit in each entry indicates the run of zeros encoded by a particular Huffman code, and the second digit indicates the size (number of bits in) a number following the Huffman code for specifying the actual value of the coefficient. For example, a run of three zeros followed by a coefficient value of 1 (a run/size combination of 3/1 in Table 4) is represented by the hexadecimal value 31 in the above array. A table similar to Table 4 may be constructed for the chrominance channel AC coefficients of the image, using different generating values, and resulting in different Huffman codes for the run/size combinations. MPEG encodes the AC coefficients differently. Rather than assign Huffman codes to run/size combinations, MPEG assigns Huffman codes to common run/value combinations. That is, common combinations of the number of zeros preceding a non-zero coefficient and the actual value of the coefficient (not just its relative size) are assigned Huffman codes. There are a very large number of possible run/value combinations, so only the most common few dozen are assigned Huffman codes. A special escape sequence handles the occasional combination that is not in the default table. The table of MPEG Huffman codes for various run/value combinations is abbreviated below: TABLE 6 Run/value Code 0/1 11 s 0/2 0100 s 0/3 00101 s 0/4 0000110 s 0/5 00100110 s 0/6 00100001 s . . . 1/1 011 s 1/2 000110 s 1/3 00100101 s . . . 2/1 0101 s 2/2 0000100 s . . . 3/1 00111 s 3/2 00100100 s . . . End of block 10 The last bit of each code, indicated by “s”, is a sign bit, with 0 indicating a positive value and 1 indicating a negative value. Combining tables 1, 2, and 6 above, it is now possible to determine the MPEG bit pattern for the luminance values of the entire example pixel block: TABLE 7 (MPEG) Coefficient Number Run/size Value Bit pattern 0 (DC) N/A −4 101 011 1 0 −5 001001101 2 0 −4 00001101 3 0 1 110 4 0 1 110 5 0 2 01000 9 3 −1 001111 End of Block 10 Clearly there is much commonality between making a JPEG image and making an MPEG I-frame. It is possible to create MPEG I-frames by creating JPEG images using dedicated circuitry in a camera, parsing the Huffman stream, and substituting the corresponding MPEG bit patterns. However, because the Huffman codes representing different DCT coefficients typically vary in length, the process of parsing the stream may be time consuming and inefficient when performed by a camera's microprocessor. The dedicated JPEG circuitry or other engine in a camera typically does not allow the compression process to be interrupted before the Huffman coding of the AC coefficients so that a different coding method could be used for construction MPEG I-frames. MPEG compression may be done without the aid of compression circuitry by a program running on a microprocessor that is part of a camera, but this method may be so time consuming that the camera user is dissatisfied. Dedicated circuitry could perform the MPEG compression quickly, but many cameras do not contain circuitry for constructing MPEG sequences, and such circuitry may be expensive. There is a need for a method of using the JPEG circuitry or other engine in a camera or other imaging device to assist in the construction of an MPEG sequence by performing the processing steps common to both JPEG and MPEG, while allowing the remaining processing to be performed efficiently. SUMMARY OF THE INVENTION A camera or other imaging device or other system configures its JPEG engine to produce a JPEG image that is in compliance with the JPEG specification but specially constructed. The configuration is chosen such that the JPEG image information for pixels of an image is stored in 8- or 16-bit groups, unlike a typical JPEG image in which image information is stored in groups of varying numbers of bits. A final software step reads the JPEG image information and constructs an equivalent MPEG I-frame. The 8- and 16-bit grouping in the JPEG image facilitates efficient conversion from JPEG to MPEG. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is a block diagram of a digital camera. FIG. 2 is a flow diagram showing the steps used to implement an embodiment of the invention. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENT FIG. 1 shows a block diagram of a digital camera. The lens ( 101 ) gathers light from a scene (not shown). The gathered light is redirected ( 102 ) to form an image of the scene on an electronic array light sensor ( 103 ). The operation of a focusing mechanism, which may include all or part of the lens ( 101 ), may be controlled by control signals ( 113 ) from a logic unit ( 110 ), which may contain a microprocessor system ( 116 ). Feedback signals ( 114 ) indicating the position of the focusing mechanism may flow from the lens ( 101 ) to the logic unit ( 110 ). A flash, or strobe ( 106 ) may be utilized to supply additional light ( 107 ) to the scene. The strobe is operated by the strobe electronics ( 108 ), which in turn are controlled by the logic unit ( 110 ). The camera may comprise a display ( 109 ) on which image data or status information may be shown. The camera may comprise a memory ( 111 ) for storage and recall of image data, as well as data interchange with other devices (not shown). The operation of the electronic array light sensor ( 103 ) may be controlled by control signals ( 105 ) from logic unit ( 110 ), and image information signals ( 104 ) flow from the sensor to the logic unit ( 110 ). Logic unit ( 110 ) may also include dedicated circuitry ( 115 ) for performing JPEG image compression. The user of the camera may operate various control inputs ( 112 ) in order to affect the operation of the camera. The camera may also comprise other controls and features that are omitted here for clarity. Tables 2 and 4 used in the example above are example tables chosen to provide good compression of photographic data. The JPEG specification allows a programmer to specify these tables for each image being compressed. This configurability is provided so that a programmer may select different Huffman codes to optimally compress data with particular statistical characteristics. In an example embodiment of the present invention, different tables are specified, but not for the usual purpose of optimal compression. The configurability of the compression engine is exploited for a different purpose than for which it was intended. For encoding the DC coefficient of each data block, the following table (analogous to Table 2) is used: TABLE 8 Y code size magnitude range 0 7 −127 . . . −64, 64 . . . 127 10 6 −63 . . . −32, 32 . . . 63 110 5 −31 . . . −16, 16 . . . 31 1110 4 −15 . . . −8, 8 . . . 15 11110 3 −7 . . . −4, 4 . . . 7 111110 2 −3 . . . −2, 2 . . . 3 1111110 1 −1, 1 11111110 0 0 11111111 8 −255 . . . −128, 128 . . . 255 Using this table of Huffman codes allows any DC coefficient between −255 and 255 to be represented with an 8- or 16-bit pattern. For example, a DC coefficient of −4, as in the above example block, is represented by the 8-bit pattern 11110 011. A DC coefficient of −1 would be represented by the 8-bit pattern 1111110 0, and a DC coefficient of 129 would be represented by the 16-bit pattern 11111111 10000001. The arrays needed to generate Table 8 algorithmically according to the JPEG specification are: Code length counts: 1, 1, 1, 1, 1, 1, 1, 2, 0, 0, 0, 0, 0, 0, 0, 0 Ordering values: 0, 1, 2, 3, 4, 5, 6, 7, 8 A similar table (analogous to table 4) may be generated for coding the AC coefficients, and is used in an embodiment of the present invention. An abbreviated version is: TABLE 9 Run/size Code 0/0 0001111111100000 (Special end-of-block character) 0/1 000111111100000 0/2 00011111100000 0/3 0001111100000 0/4 000111100000 0/5 00011100000 . . . 1/1 000111111100001 1/2 00011111100001 1/3 0001111100001 . . . 2/1 000111111100010 2/2 00011111100010 . . . 3/1 000111111100011 3/2 00011111100011 . . . Each entry in table 9 has the property that the length of the Huffman code and the length of the following bits combine to a 16-bit value to represent any particular run/value combination. Combining tables 1, 2, 3, and 9 above, it is now possible to determine the JPEG bit pattern for the luminance values of the entire example pixel block, constructed in accordance with an example embodiment of the invention: TABLE 10 (Special JPEG) Coefficient Number Run/size Value Bit pattern 0 (DC) N/A −4 11110 011 1 0/3 −5 0001111100000 010 2 0/3 −4 0001111100000 011 3 0/1 1 000111111100001 1 4 0/1 1 000111111100001 1 5 0/2 2 00011111100000 10 9 3/1 −1 000111111100011 0 End of Block 0001111111100000 The arrays needed to generate Table 9 are: Code length counts: 0, 0, 0, 0, 0, 0, 0, 16,16,16,16,16,16,16,16, 2 Ordering values (in hexadecimal notation):  08, 18, 28, 38, 48, 58, 68, 78, 88, 98, A8, B8, C8, D8, E8, F8 07, 17, 27, 37, 47, 57, 67, 77, 87, 97, A7, B7, C7, D7, E7, F7 06, 16, 26, 36, 46, 56, 66, 76, 86, 96, A6, B6, C6, D6, E6, F6 05, 15, 25, 35, 45, 55, 65, 75, 85, 95, A5, B5, C5, D5, E5, F5 04, 14, 24, 34, 44, 54, 64, 74, 84, 94, A4, B4, C4, D4, E4, F4 03, 13, 23, 33, 43, 53, 63, 73, 83, 93, A4, B3, C3, D3, E3, F3 02, 12, 22, 32, 42, 52, 62, 72, 82, 92, A2, B2, C2, D2, E2, F2 01, 11, 21, 31, 41, 51, 61, 71, 81, 91, A1, B1, C1, D1, E1, F1 00, F0 A JPEG file constructed in accordance with an example embodiment of the invention and exemplified in Table 10 has the feature that all bit patterns encoding coefficients are 8 or 16 bits in length. The data are said to be “byte-aligned”. These lengths match with the usual structure of computer memory and microprocessor architectures such that a microprocessor may read the contents of the file very quickly without the need to do bit-level manipulation of data to extract codes of varying length as would be necessary with a typical JPEG file as exemplified in Table 5. Tables can be constructed to convert a byte-aligned JPEG file to an MPEG file. For example, Tables 2 and 8 above may be combined as follows: TABLE 11 DC Coefficient Equivalent MPEG Byte value MPEG code code represented Code length 0 0000000 −127 111110 0000000 13 0 0000001 −126 111110 0000001 13 0 0000010 −125 111110 0000010 13 . . . 0 1111111 127 111110 1111111 13 10 000000 −63 11110 000000 11 10 000001 −62 11110 000001 11 . . . 10 111111 63 11110 111111 11 110 00000 −31 1110 00000 9 110 00001 −30 1110 00001 9 . . . 110 11111 31 1110 11111 9 1110 0000 −15 110 0000 7 1110 0001 −14 110 0001 7 . . . 1110 1111 15 110 1111 7 11110 000 −7 101 000 6 11110 001 −6 101 001 6 . . . 11110 111 7 101 111 6 111110 00 −3 01 00 4 111110 01 −2 01 01 4 111110 10 2 01 10 4 111110 11 3 01 11 4 1111110 0 −1 00 0 3 1111110 1 1 00 1 3 11111110 0 100 3 11111111 00000000 −255 1111110 00000000 15 11111111 00000001 −254 1111110 00000001 15 . . . 11111111 01111111 −128 1111110 01111111 15 11111111 10000000 128 1111110 10000000 15 11111111 10000001 129 1111110 10000001 15 . . . 11111111 11111110 254 1111110 11111110 15 11111111 11111111 255 1111110 11111111 15 Using a table look-up method, a program running on a camera's microprocessor system can rapidly convert each specially coded DC JPEG coefficient code to an MPEG DC coefficient code, which may be placed in a destination MPEG image file. A table similar to Table 11 may be constructed for the chrominance channels of the image. Similarly, Tables 2, 6, and 9 may be combined to form a table analogous to Table 11, but for the AC luminance coefficients. An abbreviated sample is as follows: TABLE 12 MPEG Run/ AC Code with Run/value Equivalent Code size extra bits Represented MPEG Code Length 0/4 000111100000 0000 0/−15 0000 0000 1011 11 14 . . . 0/3 0001111100000 000 0/−7 0000 0010 101 11 0/3 0001111100000 001 0/−6 0010 0001 1 9 0/3 0001111100000 110 0/6 0010 0001 0 9 . . . 0/3 0001111100000 111 0/7 0000 0010 100 11 1/3 0001111100001 000 1/−7 0000 0000 1010 11 14 1/3 0001111100001 001 1/−6 0000 0000 1011 01 14 1/3 0001111100001 110 1/6 0000 0000 1011 00 14 . . . 1/3 0001111100001 111 1/7 0000 0000 1010 10 14 . . . 0/2 00011111100000 00 0/−3 0010 11 6 0/2 00011111100000 01 0/−2 0100 1 5 0/2 00011111100000 10 0/2 0100 0 5 0/2 00011111100000 11 0/3 0010 10 6 1/2 00011111100001 00 1/−3 0010 0101 1 9 1/2 00011111100001 01 1/−2 0001 101 7 1/2 00011111100001 10 1/2 0001 100 7 1/2 00011111100001 11 1/3 0010 0101 0 9 2/2 00011111100010 00 2/−3 0000 0010 111 11 2/2 00011111100010 01 2/−2 0000 1001 8 2/2 00011111100010 10 2/2 0000 1000 8 2/2 00011111100010 11 2/3 0000 0010 110 11 3/2 00011111100011 00 3/−3 0000 0001 1100 1 13 3/2 00011111100011 01 3/−2 0010 0100 1 9 . . . 0/1 000111111100000 0 0/−1 111 3 0/1 000111111100000 1 0/1 110 3 1/1 000111111100001 0 1/−1 0111 4 1/1 000111111100001 1 1/1 0110 4 2/1 000111111100010 0 2/−1 0101 1 5 2/1 000111111100010 1 2/1 0101 0 5 3/1 000111111100011 0 3/−1 0011 11 6 3/1 000111111100011 1 3/1 0011 10 6 0/0 0001111111100000 EOB 10 2 Using a table look-up method, a program running on a camera's microprocessor system ( 116 ) may rapidly convert each specially coded AC JPEG coefficient code to an MPEG AC coefficient code, which may be placed in a destination MPEG image file. The program may also supply the proper header information to the file. FIG. 2 is a flow diagram showing the steps used to implement an embodiment of the invention. In step 200 , table generating values are chosen that will generate Huffman code tables, such as Tables 8 and 9 above, having the property that the Huffman codes for each run/size combination and the additional bits to fully specify the value of each coefficient combine to form a code that is exactly 8 or 16 bits long. In step 202 , the table generating values are provided to the JPEG engine to be used in performing JPEG processing. This configures the JPEG engine. In step 204 , an uncompressed image is obtained. The image may be obtained by taking a photograph with a digital camera, making an image with some other imaging device such as a scanner, or even reading a previously stored digital file. In step 206 , the JPEG engine performs the JPEG processing, generating the code tables from the table generating values supplied earlier and performing the steps of the JPEG technique. This creates a “byte-aligned” JPEG data stream. In step 208 , the byte-aligned JPEG data stream is read by a program running on a processor. In step 210 , the byte-aligned data stream is efficiently converted by a simple table lookup or similar means to an MPEG data stream. In step 212 , the MPEG data stream is stored as an MPEG I-frame. This step may also include adding header information to the file. The foregoing description of the present invention has been presented for purposes of illustration and description. It is not intended to be exhaustive or to limit the invention to the precise form disclosed, and other modifications and variations may be possible in light of the above teachings. For example, the invention may be used to take advantage of existing JPEG software code libraries for aiding in the construction of MPEG I-frames. No hardware need be involved at all; the JPEG engine may be entirely software based. The embodiment was chosen and described in order to best explain the principles of the invention and its practical application to thereby enable others skilled in the art to best utilize the invention in various embodiments and various modifications as are suited to the particular use contemplated. It is intended that the appended claims be construed to include other alternative embodiments of the invention except insofar as limited by the prior art.

A camera or other imaging device or other system configures its JPEG engine to produce a JPEG image that is in compliance with the JPEG specification but specially constructed. The configuration is chosen such that the JPEG image information for pixels of an image is stored in 8- or 16-bit groups, unlike a typical JPEG image in which image information is stored in groups of varying numbers of bits. A final software step reads the JPEG image information and constructs an equivalent MPEG I-frame. The 8- and 16-bit grouping in the JPEG image facilitates efficient conversion from JPEG to MPEG.

CROSS-REFERENCE TO RELATED INVENTIONS This application claims the benefit of provisional application serial No. 60/080,790, entitled “Application Broker” filed Apr. 6, 1998. Pending this application is a continuation-in-part of the application Ser. No. 09/190,757, now U.S. Pat. No. 6,104,392 entitled “Method of Displaying Applications on a Variety of Client Devices in a Client/Server Network,” filed Nov. 12, 1998. STATEMENT REGARDING FEDERALLY FUNDED RESEARCH Not Applicable. REFERENCE TO A MICROFICHE INDEX Not Applicable. COPYRIGHT NOTICE Copyright 1999 The Santa Cruz Operation, Inc. A portion of the disclosure of this patent document contains materials that are subject to copyright protection. The owner has no objection to the facsimile reproduction by anyone of the patent document or patent disclosure, as it appears in the Patent and Trademark Office patent file or records, but otherwise reserves all rights, copyright rights whatsoever. BACKGROUND OF THE INVENTION 1. Field of the Invention The present invention relates to networked data processing environments using a client/server architecture, and, in particular, to client-server systems where there exist one or more clients of varying capability connected via network connections of varying bandwidth and latentcy to one or more servers providing application program services or database services to the connected clients. 2. Background Information Until the 1980s the computer network typically was used as a means of connecting to a large mainframe environment using dedicated hardware terminals and proprietary protocols. Next UNIX servers grew in usage and with it came standardization of the networking protocols, in particular TCP/IP (Transmission Control Protocol/Internet Protocol). Simultaneously, there was a shift in the computing paradigm towards client/server architectures. This allowed the processing power to be distributed over the network and not be limited to servers which could not scale to meet the growing number of users and their increasing demands. This lead to the need for clients to become more intelligent and powerful. Significant desktop clients came into use: Microsoft Windows. Each iteration of this operating system brought about more functionality requiring more powerful desktop clients. More and more software had to be installed on these clients leading to what is termed “fat clients” and each client required individual configuration. People began to find that the amount of time and money required to maintain these powerful clients was increasing. High availability networks together with hardware and software needed to support such networks have become the norm. In these environments there is not a homogeneous structure of one type of server and one type of client, but a variety of such devices. Within a network there is a wide variety or servers and clients. Upgrades and new applications for this diverse mix of clients usually require that each client be individually upgraded. Also each user has specific needs on the desktop client: configuration, security, access control, mobility. The information services department provides this by administering each client separately. If remote offices are involved, costs to do this increase significantly. Performance in this heterogeneous client network must also be maintained. Slow performance takes up time and money. Networks vary in bandwidth, e.g. modem links, ATMS, Frame Relays, etc. The World Wide Web (the “Web”) has come to the forefront in the current era of the Internet/Intranet and networks have become an integral part of day to day work. Modem speeds double every year and 100 Mbit/sec Local Area Networks have arrived. The Web is now a well-accepted medium for publishing information, in the form of text and graphics (including sound and video). Web programming languages such as Java, JavaScript, and CGI have now extended the Web to applications. This is fine for new applications but existing applications also need a route into this world. Existing applications have either had a character-based or windows-based user interface. Now such applications need a web user interface. The web user interface provides a presentation layer to the user of the application. It must provide an input/output method for the user to interact with the application. There are a number of ways to do this including: (1) HTML (HyperText Markup Language) replacement of the current user interface; (2) Non-Java plug-ins; and (3) Java-based emulation. The first solution involves rewriting the application. The second involves installing more software on the clients leading to “fat clients.” The third is preferred. A large number of vendors offer a character or graphical emulation package that runs on desktop clients such as Windows, UNIX, etc. These emulators could be rewritten in Java and such Java emulator will run on just about any client. However, this approach leads to fat Java clients. For performance reasons, users will not want to wait for large Java applets to download. These Java applets will grow in size, as users demand more and more functionality. Storing these Java applets locally solves the performance problem but leads to fat clients. If state information is stored on the client leading each client to have its own particular configuration parameters, the problem of fat clients is further exacerbated and each client is being managed individually rather than from a central place on the server. Web browsers have an API (Application Programmers Interface) enabling software developers to provide helper applications that allow users to run applications or view unsupported document types on their client platform. These are termed “browser plug-ins.” They are both platform-specific and browser-specific. For example, for two platforms (Windows and UNIX) and two browser types (Netscape and Microsoft), four different implementations of the plug-in would be needed. This type of solution is not cross platform and the majority of these plug-ins are proprietary (e.g., Microsoft ActiveX) locking users into vendor specific solutions. Once the web display of an application is possible, the next step is to make it available to all users. A number of methods that could be used include: (1) local installation; (2) push technology; (3) on-demand access. Local installation involves an administrator installing the application or connectivity software on every single client. This is disadvantageous in that it makes the clients more difficult and costly to manage and leads to fat clients. Push technology involves storing all the files and data associated with users and applications on a server and transmitting them out via virtual channels on a network to clients. Where all storage is on the sever, users experience poor performance in waiting for applets and applications to execute, or downloads are cumbersome and in some cases unusable. Local storage of applications and state information is used to improve performanc. This approach starts off well by using a central server but when applications and any associated state are stored on the client, the fat client problem arises. With on-demand access as used in the present invention, user state, applications, connectivity software and the associated configuration data are stored on a server. Applets are downloaded on demand when the clients request an action, such as start up an application. All state information is kept on the server and can then be managed centrally rather than individually on each client. Keeping state information on the server also makes the client resilient. If the client connection is lost or the client itself is replaced, nothing is lost and no replication is needed. Next the applications and data must be made available to selected users in a secure manner. For manageability, this is done centrally on a server and made available by the most common medium to all users, the Web. But this raises more questions: What do web pages on that server have to contain? What editor has to be used? How do you make it available to selected users or groups of users? Are there different web pages for different types of users? Where does the user profile and application configuration reside? Do users authenticate themselves every time they want to run applications? How is the authentication done? What if the user is already authenticated on the server and does not want to do again and again for each application? What if all of this information is already available and duplication is not desired? As the above list shows, providing a display mechanism via Java emulation is only a partial solution to web-based delivery of applications. To achieve optimal performance on all networks is difficult. Protocols are designed with specific functions in mind, e.g., inter-process communication, graphics rendering, etc. Protocols are rarely designed with the goal of providing uniform performance over complex network routes that have different permutations of bandwidths. To choose the right protocol a number of assumptions could be made. For a high bandwidth network, such as a fast LAN with low latency, the X protocol works well but is unusable over a slow modem link. For a low bandwidth network, such as a slow WAN or modem connection, compression can be used to optimize performance. The ICA protocol of Citrix Systems, Inc. works well over a slow modem connection but is inefficient on a fast LAN connection. It would be advantageous to be able to deliver upgrades to existing applications or roll out new applications to diverse clients; to be able to centralize the administration of clients and their applications; to be able to balance loads in a heterogeneous client server network and to provide high security and manage sessions across heterogeneous servers and services. SUMMARY OF THE INVENTION The inventive universal application server, also termed an application broker, is added in a nondisruptive manner to an existing computer network environment forming a second server tier intermediate the application servers tier on a third tier and the client devices on a first tier. The system provides the following elements: a shared database that is used to describe the location of the application programs, the application servers, the users of the system and a description of which users are provided with access to which application programs; a set of protocol engines processes which support industry standard display protocols (X11, telnet, VT220, etc.) which form the protocols as used by the application programs; a set of display engines processes which are designed to operate within the various client device environments and which communicate with the PE to provide presentation and user interaction services with the users of the client devices; a set of delivery servers which can be used to deliver the display engines to the client devices on demand; and a set of management engines—a data store engine, a status engine, a session manager engine, an administrative engine that allow manipulation of the shared database and control of the operation of the system. Once the application broker has been installed into the existing client server environment it operates in the following manner. Management processes define and maintain the shared database. This consists of making entries for each user of the system for each application server and for each application program. Users connect to the application broker from their client devices and are authenticated against the database. Display engine components are automatically downloaded to the client device, and are used to present the user with a webtop (visual menu) of applications that are available for execution. The contents of the webtop are specific to the user and are constructed dynamically from the database. The user chooses which applications to run. This information is passed from the display engine running on the client device to the universal application server The universal application server, using information stored in the database determines the location of the application program and communicates using industry standard protocols with the appropriate application server. This communication causes the launch of an instance of the application program to take place. The application program launches and is instructed to perform user interaction over a standard protocol by the universal application server. The application broker also arranges to create an instance of the appropriate protocol engine to handle the specific protocol used for a number of applications, such as X-windows applications. It also downloads a suitable display engine to the client device to interact with the user. The protocol engine takes output requests from the application program and converts them into a form that is suitable for use by the downloaded display engine. This information is then forwarded to the display engine where it then converts the information into a form suitable for display within the client device environment. The user interactions with the display engine and the result of these actions (typing, mouse manipulation etc.) are transmitted to the protocol engine. The protocol engine converts these user inputs into the standard protocol used and passes them to the application program. As set forth in the copending application referenced above the protocol engines and display engine processes transmit to and display on the display device of the client device, display requests of the requested application service which are not supported by the client device. The protocol engine is initialized with parameters of the client device comprising display operations supported by the client device, the relative cost of each supported operation and a local performance factor and parameters of the connection from the server to the display engine comprising bandwidth and latency. The protocol engine maintains a first queue for retaining pending display requests and a second queue for retaining transmitted display requests of the requested application service for display on the client device. Each request has a corresponding relative cost and request time determined by the relative cost and the local performance factor. The protocol engine also maintains a total request time for all requests in both queues and maintains a total network time for all requests in the second queue. While the total request time is less than a predetermined first value, the protocol engine accepts new pending requests from the application service that are to be displayed on the client device. The protocol engine executes the new pending display request to create an image to be displayed on the client device, saves the resulting image to memory and determines for the image its relative cost and its corresponding request time based the local performance factor received from the display engine. The protocol engine also converts the new pending display request from the application service into a sequence of converted requests that are supported by the client device, places the sequence of converted requests in the first queue, determines for each converted display request in the sequence its relative cost and its corresponding request time based on the local performance factor, and increments the total request time by the request time of each added request. Further while the total network time is less than a predetermined second value optimizing the first queue using a merge optimization occurs when the request times of the pending requests in the first queue exceed a predetermined third value. The total request time is updated based on the results of the merge optimization. Next a pending request is read from the first queue and labeled with a first sequence identifier. Each read request is then encoded for transmission to the display engine over a network connection. A copy of the optimized request is placed into the second queue and the total network time is incremented by the request time for the newly added request. The encoded display request is transmitted to the display engine. The display engine unencodes the received converted display request and displays it on the client device. Next it generates a second sequence identifier corresponding to first sequence identifier of the displayed request; and periodically transmits to the protocol engine the second sequence identifier of the last received display request displayed on the client device. Upon receipt of the second sequence identifier, the protocol engine further comprises the steps of: deleting from the second queue the display request whose first sequence identifier is the same as the second sequence identifier and all pending displayed requests in the second queue having first sequence identifiers that are prior to such second sequence identifier; and decrementing from the total request time and from the total network time the request time values of each of the deleted display requests. The application broker using the information contained in the data store regarding the application programs and the application servers used to run them is able to keep track of which third tier server is actually running the application program. The number of all applications running on each application server is tracked. At launch time of the application program, the application broker chooses the application server from the list associated with the application that currently has the lowest number of applications running on it. This selection process provides a load balancing mechanism for the application servers. BRIEF DESCRIPTION OF THE SEVERAL VIEWS OF THE DRAWINGS For a better understanding of the invention reference may be made to the preferred embodiments of the invention shown in the accompanying drawings where: FIG. 1 illustrates a Universal Application server (the “UAP”) system comprised of a client/server network showing the interconnection of the UAP server to other applications and database servers and to a diverse set of clients. FIG. 2 is an illustration of the Adaptive Internet Protocol (AIP) Link as implemented in the UAP system. FIG. 3 shows the processes operating in the display engine and protocol engine components of the present invention. FIG. 4 is a more detailed data flow diagram of the AIP link of the present invention; FIG. 5 is a table of sample of AIP requests. FIG. 6 presents a table of AIP request translations and a table of relative costs for performing an AIP request. FIGS. 7A and 7B illustrate a flow diagram of the startup process of the application broker system. DETAILED DESCRIPTION OF THE INVENTION Overview of a Universal Application Server System As used in this application the term “engine” is used to refer to a process, method or series of related processes or methods for performing certain actions. For example the session engine is a process that is used to control a current session between the server and a client device. Also the terms “universal application server,” “UAP server,” and “application broker” are used interchangeably. Also the same or similar item occurring in more than one figure carry the same or similar numeric designation. As shown in FIG. 1, the network 10 comprises three tiers used to illustrate the connections of the various inventive processes and routing of data information among the various processes and components comprising the network. First tier 12 contains a variety of diverse client devices, generally indicated at 14 , having different interfaces, such as PC 14 A, UNIX computer 14 B, network computer 14 C, all illustrated as having a Java Virtual machine interface 16 to second tier 40 , or a thin client 14 D shown having a browser interface 18 . The inventive processes are also useable in a standard Microsoft Windows environment that uses interfaces such as Win16 and Win32. In addition, administration computer 20 , illustrated as a network computer, is also grouped in first tier 12 and is in communication with processes running on second tier 40 . Second tier 40 comprises UAP server 50 having various engines or processes executing thereon together with various interfaces to first tier 12 and to third tier 80 . The interfaces to first tier 12 include X windows interface 52 , character based interface 54 , Java interface 56 , and HTML interface 58 . The interfaces to third tier 80 include an X Windows interface 60 , terminal emulators 62 , web server 64 , and lightweight directory access protocol (LDAP) interface 66 . The third tier 80 includes the various application servers, generally indicated at 82 , including UNIX X Windows server 82 A, Windows NT server 82 B, UNIX and IBM character-based application servers 82 C, directory services servers 82 D, the Internet 82 E and various data base servers 82 F. UNIX X window server 82 A and the Windows NT server 82 B interface to second tier 40 via X Windows interface 60 . UNIX and IBM character-based application servers 82 C interface via appropriate terminal emulators 62 . Web server 64 provides the interface to the Internet 82 E while the directory services server 82 D uses the LDAP interface 66 . Although all of the application servers are shown in FIG. 1 as being in a single tier, it should be realized that the application servers can also unction as clients to other servers not shown and that the requested application may reside in a server located elsewhere in the network. Also shown are session manager engine 70 that controls active sessions between client devices 14 and application servers 82 and the administrative engine 72 which controls the operation and administration of UAP server 50 and which is in communication with the administrative computer 20 operated by a system administrator. Client connections and requests, 30 , are routed to the appropriate interfaces on UAP server 50 which in turn processes them using one of more of the UAP engines described below and obtains the requested service or data from the appropriate application server 82 . UAP server 50 then returns to the client devices 14 display requests from the requested application and other data. UAP server 50 resides on a host or server on the network that has a web server running on it. It can be viewed as a black box sitting somewhere on the network that enables any client to access any service by providing the intelligence needed to do this. A variety of naming standards and methods are used to access resources on the network. Computers on the Internet use conventions such as Domain Name Service (DNS) or Windows Internet Naming Service (WINS). File naming conventions are tied to operating systems of which there are many. For example, Microsoft operating systems use the Universal Naming Convention (UNC). These resources are disparate and different entry points and methods are needed to access them. The UAP server provides an integration point for a wide variety of naming standards by federating the name spaces. It provides a single point of entry to access any network resource. The X/Open Federated Naming (XFN) based scheme is the preferred naming standard to permit such integration amongst these variety of naming services and standards. Large and complex information about network resources is usually kept by directory servers, for example, Novell® NDS™, Microsoft Active Directory, Netscape Directory Server™. UAP server 50 does not provide a complete replacement or duplication of these services. UAP server 50 provides a thin layer that interfaces into these services using LDAP (Lightweight Directory Access Protocol). LDAP is a current defacto industry standard interface to directory services. By doing this UAP server 50 avoids duplication of information and functionality which may already be present on the network. It would be unnecessary and time-consuming process to replicate this information. For example, if a thousand users are already set up, the use of the LDAP layer allows use of this information by UAP server 50 and administrative engine 72 , rather than making administrators recreate it and double the maintenance load. As illustrated in FIG. 2, UAP server 250 contains the following processes: administrative engine 272 ; status manager engine 274 , protocol engines, 276 A- 276 C, display engines 278 A- 278 C; web server 264 and webtop storage 268 for the storage of user webtops; data store 273 containing various user and system data such as the LDAP protocol, user passwords, user application objects, startup engine, and other data related to the operation of UAP server 50 , data store engine 271 and one or more session managers 270 , which can be configured to permit suspendable and resumable sessions. Login port 251 , reconnect port 253 , a plurality of session ports 255 are provided for communication with client devices. Display engines, generally designated 278 , are stored on UAP server 250 and are downloaded to client device 214 when needed. Client device 214 is a computer system that has a display and input devices such as a keyboard, mouse, touchpad, etc., as is known in the art. Shown in the third tier 280 are three application servers, generally designated 282 , X Windows server 282 A, Windows NT server 282 B, and character-based server 282 C. UAP server 250 provides administrative engine 272 . Only authorized administrative users have the ability to run this process. The main functions of the administrative engine 272 are: Publishing applications and documents to users Organizing user profiles (what each user is allowed to access) Viewing and changing the contents of the organizational hierarchy Monitoring and controlling sessions, including viewing which users are logged on and what applications are running, and stopping applications Configuring the UAP server settings such as logging and diagnostic specification, authentication mechanisms and file locations Sitting at the heart of UAP server 250 is a suite of server engines or processes that coordinate all functions. These are status manager engine 274 , data store engine 271 , session manager engine 270 , protocol engines, generally designated 276 , and the display engines 278 . These engines cooperate to provide the central point of access to the client devices and handle the following: Security management Retrieval and storage of state associated with users and applications Invocation of applications and their supporting display mechanisms Session management Status manager engine 274 provides initialization and control over all the other server engines. It is the first process that is launched when UAP server 250 starts up. It is designed to manage all other processes and act as a central point of brokering requests and actions. Data store engine 271 provides the initial interface to a connecting client device. It controls the contents of data store 273 , which holds information about services, network devices, as well as users and their associated applications and data. UAP server 250 uses a bootstrap applet 221 that downloads itself to the client device 214 on demand from a user. This applet then connects to data store engine 271 , which provides an authentication service, after which requests can be issued to the data store engine 271 to: generate webtops; configure the webtop; invoke applications; retrieve documents; and, if requested by the system administrator, view and modify the contents of the data store itself for system administration. For each application, data store 273 contains a list of the application servers that are able to run this particular application. The data store is a hierarchy based upon the X 500 standard. The major items contained with the database are: People. These represent the users of the system. The location of the person within the hierarchy can be used to reflect the organization of the company. In particular, it is possible to group people to provide for easy administration of people that have access to similar application programs. Information stored for a person includes authentication information such as application passwords and login password, a list of groups to which that person belongs and a list of applications that should be presented to that user. Applications. These items represent the application programs that are deliverable by the system. The location of an application program within the hierarchy is used to determine which users will be presented with that application. Information present for an application includes, protocols used by the application (X11, Telnet, etc.), location of the application (application server name, location within the file system on that application server) and configuration information required to launch the application (environment settings for the application server, command line arguments, display size information etc.). Groups. These items are used to form convenient administration objects. They can be used to simplify the delivery of applications to sub-sets of users or to provide access of sets of applications. The actual database itself is federated and can access data from several sources. These include directory services (via LDAP protocols), local system authentication information (passwords etc.), and dynamic state information (running applications etc.). At any given time while UAP 250 server is active, session manager engine 270 contains dynamic information about which users are connected and the applications they are running. In addition, session manager engine 270 is responsible for invoking applications after the client device request has been processed by data store engine 271 . Session manager engine 270 can be viewed as the keeper of all live information in the UAP server 250 . Session manager engine 270 tracks which applications are running on which application servers. At launch time of a given application, session manager 270 obtains the list of application servers capable of running this application and chooses the application server that can run the requested application and that currently has the lowest total number of applications running. This selection process provides a load balancing mechanism. More complex load balancing schemes can be used but we have found that this process produces a reasonably even sharing of the processing load amongst the application servers. Protocol engines 276 A- 276 C and display engines 278 A- 278 C provide the emulation necessary for the user to view and interact with applications. Preferably, UAP server 250 via session manager 270 associates one protocol engine and one display engine for each corresponding application type. They are invoked and the display engine is downloaded on-demand to the requesting client device. Protocol engines 276 A- 276 C run on UAP server 250 and perform the bulk of the emulation by acting as a client of its associated application type running on the network. The protocol engine understands the standard protocols the application types currently use. X-windows type applications would use a protocol engine that is different than the protocol engine used for Microsoft Windows type applications. The appropriate protocol engine translates the standard protocol requests into an adaptive Internet protocol that the display engine on the client device can display. Using this method existing applications continue to run on servers they are currently installed on, untouched and without any re-engineering or re-writes needed in order to function with the client device. Preferably, the protocol engines are implemented as native binaries to ensure optimal performance on UAP server 250 . Display engines 278 A- 278 C are preferably Java applets that are downloaded on demand by the client device. The display engines are small in size—around 200 Kilobytes—and render the application on the client device display and allow input/output to the user. As a result the display engines are quick to download, even over low-bandwidth networks. The display engines know nothing about the state of the application that is being run. There is a display model mismatch between the existing application environments and the display model used on the display engine. Operations that are not supported on the display engine (e.g., plane-masking, logical operators etc.) are carried out on the protocol engine and sent down as supportable operations (e.g., bitmap copy from memory) to the display engine. The most common mode of access to the UAP server is from web browsers. Users of UAP server 250 interact with applications and documents on the network using the web equivalent of a desktop—the webtop. UAP server 250 gathers all objects (applications, documents, etc) associated with a user and dynamically creates a web page to represent this information. This web page contains smart applets, preferably written in the Java programming language, represented as graphical icons (see 225 A- 225 C). When the user clicks on these smart icons, requests are issued to invoke services or applications or view documents. The navigation model that is used in browsers involves going back and forward between pages. From time to time pages with applications running inside them may be out of view on the display screen of the client device. Browsers temporarily cache Java applets such as the display engine inside these pages that are out of view. After a certain amount of time these applets are cached out by the browser. UAP server 250 works around this problem by providing an automatic reconnect of the display engine without shutting down the application (as indicated at path P 10 between display engines 278 A-C and reconnect port 253 ). UAP server 250 uses session manger engine 270 and sessions to provide and control the access and interaction of the user with the desired network application or database. A session is associated with each application or database the user may be running. An application session can be configured by administrators to be resumable. With a resumable session users can disconnect themselves from UAP server 250 but the resumable applications remain running. When the user reconnects to the session, the application is presented in the same state as it was when left. This reduces start up time for applications and allows client resilience. For example, some applications take a long time to start up. At other times, after starting the application, the user goes through a laborious route to a particular point, for example, a character-based data entry application with a large number of menu options or a processor intensive query or calculation. In a network-centric environment, the connection from client to server is critical. If the connection drops for any particular reason, for example, if a modem connection is interrupted, the server must be able to recreate the state associated with the client device when it reconnects. The suspend/resume facility of session manager 270 allows this to occur. During the configuration process, administrators can choose to disallow the resume facility, on a per application basis, to save server resources. For example, a simple calculator application is unlikely to need resuming. Administrators also have the ability to terminate application sessions while they are running. Operation of the UAP Server The process used by the UAP server to allow client device 214 to connect, deliver the client's webtop and run applications or view documents is illustrated in FIG. 2 . Client device 214 connects to web server 264 as shown at path P 1 . Here the user is shown connecting to web server 264 via a Java enabled web browser 218 but the user can also connect by other means including using a Java virtual machine running on the client device or use of a native binary application. Web server 264 retrieves from data store 268 a web page containing UAP bootstrap applet 221 and returns it to web browser 218 . Bootstrap applet 221 is executed on client device 214 . This applet connects with data store engine 271 via the login port 251 as shown in path P 2 . Data store engine 271 then provides client device 214 with login applet 223 via path P 2 . Login applet 223 is executed on client device 214 and the user is authenticated to UAP server 250 using this applet. If the correct user-name and password are supplied, applet 223 issues a request to data store engine 271 to search and find the webtop for that user. A webtop 225 is created dynamically using the objects associated with the user stored in data store 273 and is loaded into browser 218 running on client device 214 . For a new user, the system administrator provides the initial objects to be used in the construction of the web page together with the level of user privileges available to the user. Session manager 270 begins a bottom up traversal of the data store hierarchy starting from the user and expanding any groups. During this traversal a list of applications specific to the user is created. When the traversal is complete, the list of applications forms the set of applications to be made available to this user. Session manager 270 transmits presentation data for the application set to the login applet 223 (this includes textual and graphical representations, i.e., a web page or webtop) that, in turn, presents the application set in a manner appropriate to client device 214 . The user chooses an application to launch (via some sort of input device supported by client device 214 . When the user clicks one of the icons ( 225 A- 225 C and preferably a smart Java icon) webtop 225 , the web page associated with that object is downloaded from web page store 268 and web server 264 . If the object is an application, then the appropriate display engine for the client device is downloaded. In FIG. 2, three display engines 278 A, 278 B, and 278 C are shown on web browser 218 indicating that this user is using three applications. Each display engine 278 performs an initialization routine with client device 214 to determine supported display operations and display performance factors. Preferably, each display engine 278 contains a table 279 of display operations that are expected to be supported by the client device in its operating environment. The display engine performs a test of the client device to determine which of the display operations in the table can actually be performed on the client device and modifies table 279 to indicate which operations are supported on the client device 214 . The modified table 279 is stored in the display engine and is also sent to UAP server 250 for further use. The display engine also calculates a local performance factor or scaling factor for the client device to determine the relative cost of performing the supported display operations. Preferably, two bitmap copy operation tests are performed, i.e., data in memory is written to the display. In the first test, a small number of pixels, one or two, are written to from memory and the time to perform this operation is noted. This step is repeated an arbitrary number of times and an average time T small is determined. This time, in microseconds, represents the overhead to work with a pixel. In the second test a larger image involving a larger number of pixels is copied from memory to the screen and the time to perform this T large is noted. Next the Per Pixel Cost, PPC, is determined as follows:   PPC=[T large −T small ]/[Number of pixels written in second test] This local performance factor, PPC, is communicated to UAP server 250 for use by the protocol engine in determining request times for the requests found in the two queues. Each protocol engine uses database engine 271 and session manger 270 to determine which application server to contact and the protocol to use to launch the application. Each protocol engine connects with the application server identified by session manager 270 and starts the application, passing identification data to ensure that the launched application transmits user interaction and display requests to the correct instance of protocol engine. For each application the data store 273 holds a list of application servers that are able to run the application that is being requested by a particular instance of a protocol engine. This list is used by the launch process to determine which application server in the third tier is used to actually run the third tier application program. Session manager 270 keeps track of the total number of applications running on each application server. At launch time, session manager 270 selects the application server for the list associated with the application that currently has the lowest number of applications running on it. This selection process provides a load balancing mechanism for the application servers in the list so as to reasonably even the load among the application servers. Display engine 278 A issues a request to data store engine 271 to find the application, in this case, the X Windows application 282 A. This request is passed onto session manager engine 270 which checks to see if the requested application is running and needs to be resumed to the user. If an appropriate protocol engine is not already running on behalf of this user one is started by the session manager. Protocol engines are able to handle multiple applications that use the same protocols and so only one of each application type is required per user. An X Windows protocol engine can be used by multiple X Windows applications but not by a Windows application. Session manager engine 270 starts the requested application on the network using a password stored in an encrypted cache (not shown) in data store 273 if it has already been supplied. If the password has not been supplied, session manager engine 270 prompts the user via display engine 278 A for a password that is then stored in the encrypted cache. Session manager engine 270 then invokes the correct protocol engine for the requested application type, in this case protocol engine 276 A. Protocol engine 276 A establishes a new connection, as shown at path P 3 , to the requested application 282 A that had been previously started by session manager engine 270 . Session manager engine 270 stores the information about the protocol engine 276 A used with application object 282 A in data store 273 . As part of this sequence, if the application is one in which the application needs to connect to the protocol engine (which is the process used for X Windows applications) then information that is required by this mechanism is placed in the application's environment to allow it to connect to the appropriate protocol engine. Next because application object 282 A is an X Windows application type, application 282 A connects to protocol engine 276 A. Application 282 A is then in a suspended state. Information is passed back to the display engine 278 A in order to enable it to connect back to protocol engine 276 A. If the application is already running then it will be already connected to an active protocol engine and will be in a suspended state thus the startup steps to establish a connection to the application are not required. Similar actions are performed by display engines 278 B and 278 C to request use of Windows application 282 B and character-based application 282 C, respectively. Instances of protocol engines 276 B and 276 C are established by session manager 270 , and protocol engines 276 B and 276 C establish connections to applications 282 B and 282 C, respectively, along paths P 4 and P 5 respectively. For the character-based application types, the open connection to the application is passed to the protocol engine. Display engines 278 A- 278 C on client device 214 connect to ports 255 A- 255 C on the UAP server and authenticate themselves. Display engines connect back using the information sent back from the appropriate protocol engine (this includes port number and temporary authentication information). This mechanism ensures that a separate connection is used for each display engine and avoids the need to multiplex data over a single connection. Status manager engine 274 receives these requests from the ports 255 via path P 6 and passes it onto session manager engine 270 . Session manager engine 270 then connects up the requesting display engine to the appropriate protocol engine. In this case display engines 278 A, 278 B, and 278 C are connected by session manager engine 270 to protocol engines 276 A, 276 B and 276 C, respectively via paths P 7 , P 8 and P 9 , respectively, and ports 255 A, 255 B and 255 C, respectively. Login port 251 remains open and is used for later launch requests and updates to the webtop. The launch process described above is illustrated in the flow diagram presented in FIGS. 7A and 7B. There the launch process 700 starts at step 702 . The process proceeds to step 706 where status manager engine 274 is started. Next at step 710 , all other engines are initialized. The process proceeds to step 714 where the application broker or UAP server is waiting for a request from a client. If no request is received the process proceeds to step 718 to wait for a client request and then back to step 714 to see if a client request has been received. If a step 714 a client request has been received, the process proceeds to step 722 where web server 264 retrieves bootstrap applet 221 from web pages store 268 and download it to client device 214 . At step 726 the user presents login information such as user name and passwords using bootstrap applet 221 . At step 730 the user information is checked and if it is not approved, the process proceeds to step 734 where the login is rejected. Here the process loops back to step 726 to receive user information. In the event of a rejected login various known mechanisms such as allowing multiple retries or sending the login information to the system administration and management engine 72 for use by the systems administrator. If at step 730 , the login information is correct, the process proceeds to step 738 where the session manager downloads to client device 214 a webtop that is built up using the bottom up traversal of the data store 273 containing icons representing the application programs available to the user. The system at step 742 waits to receive a request from a client for an application. If none is requested the process proceeds to step 746 to wait for a client application request and then loops back to step 742 . If an application request is received at step 742 the process proceeds to step 750 where a check of the data store for the user's password is performed. If no password is found the process proceeds to step 754 where a password request screen is downloaded to client device 214 asking the user to create a password. At step 758 it is determined if the user has submitted a password. If not, the process goes to step 762 where the session is ended. If at step 758 if a password is given the process, at step 766 , encrypts it and places it in data store 273 to be used for future logins. The process proceeds to step 770 . At step 750 if the password has been found in the data store, the process proceeds to step 770 . At step 770 , a determination is made if an existing session is in place for this user and the requested application. If it is, at step 774 , the user is connected to the existing session, If not, a new session is started at step 782 and the session manager performs the load balancing previously described. Proceeding from either step 774 or 782 the session manager, at step 778 , instantiates the appropriate protocol engine if one is not already running for the requested application type, instantiates display engine for the client device type, and downloads the display engine to the client device. Also here the instantiated protocol engine establishes a connection to the application server and the requested application. The user then interacts with the requested application. A determination is made at step 790 whether or not the application is resumable in the event that the connection to the client is ended. This information is either retrieved from data store 273 and can be determined from user input gotten when the application is started. Disconnection's can be voluntary by the user or accidental due to network problems. If the application is not to be resumable, the process goes to step 794 to disconnect the application upon the user exiting. If the application is resumable, the process at step 798 will keep the connection to the application server and suspend the instance of the application that was running until the client device and user reconnect. The application program that is instantiated on the application server residing on the third tier can be a program such as, for example, a drawing program, a spreadsheet program, a database program. The nature of the application program is not critical to the operation of the universal application server. Also the connection between the protocol engine and the third application server can be extended to other servers as is known in the art to order to access the desired application program. This routing information would be stored in the datastore. The connection between the requesting display engine and the protocol engine forms an adaptive Internet protocol link. The first phase in this link is to pass parameters identifying the characteristics of the client device and network connection to the protocol engine. These include the local performance factors and the supported display operations table, and the bandwidth of the connection. The adaptive Internet protocol link then tunes itself for optimal performance for the given connection and client device. Finally, the display screens of the application are displayed on the client using display operations supported by the client device. The adaptive Internet protocol link and, in particular, the protocol engine, monitors any changes in the network connection between the protocol engine and the display engine. Adaptive Internet Protocol (AIP) Link Operation The AIP link will be described in detail in relation to protocol engine 276 A and display engine 278 A. To deliver access to multiple types of client devices over a variety of network connections, the AIP link has to adapt to the characteristics of the client device and the network it is on. The AIP link does this by providing heuristic mechanisms that optimize the responsiveness of applications by monitoring, measuring and adapting the ways in which data is transferred between applications and the client devices including: Round-trip measurement, latency and bandwidth “Just fall” transmission window management Use of machine independent drawing capabilities (e.g., Java) for specific graphics functions Color management Data compression Server-managed client caching Queue management for graphics operations Metering adjustment for “busy” networks Suspendable/resumable sessions The goal of the AIP link is to maximize use of the connection to the client device by trying to limit the number of transmitted requests while not overloading the connection and to optimize the requests to the fullest extent possible while not underutilizing the connection. FIG. 3 shows the major process structure of the AIP link 300 . The AIP link 300 comprises two major processes working cooperatively: Protocol Engine (PE) 276 A and Display Engine (DE) 278 A. The initial conditions are for PE 276 A to be running on UAP server 250 waiting for DE 278 A to contact it. The two processes are connected via a network connection using a standard network transmission protocol such as TCP/IP (path P 7 ). Following this connection, the two exchange a series of requests to establish the characteristics of the client device, e.g. type of device, operating system, and the network connection, e.g. bandwidth and latentcy. DE 278 A sends to PE 276 A modified table 279 that describes the display operations or graphics primitives supported by client device 214 and the relative cost of executing each of these operations or primitives. In the following description of the various processes the numbering of the steps is for convenience and to ease understanding of the processes. It will be apparent to those of skill in the art that the sequence of various steps may be changed without affecting the overall operation of the process. As show in FIG. 4, in PE 276 A, process 310 , at step 1 waits for the Total Request Time (TRT) for the selected application to be less than a first pre-set limit and if it is it accepts from the selected application a received display request 312 formatted in a standard network display protocol, such as the X11 Windows protocol. The TRT represents the time costs for all of the requests currently found in PE 276 A, specifically those in the Pending Request Queue (the first queue) and the Transmitted Request Queue (the second queue) described below. At step 2 , process 310 executes received application display request 312 and writes the resulting image to memory. At step 3 , using modified table 279 of supported display operations previously sent from DE 278 A, process 310 performs various transformations as is known in the art on request 312 converting it to a series of AIP requests which are simpler operations that are supported by client device 214 . FIG. 5 presents a table listing some of the AIP requests that are available for use in the conversion process. This table is not meant to be an exclusive listing of all available AIP operations and other operations can be used as will be appreciated by those of skill in the art. Steps 2 and 3 can be performed in any order. At step 4 costing is determined for the image and for each of the converted requests created in steps 2 and 3 , respectively. Process 310 utilizes the performance factor, i.e., the PPC factor, received from DE 278 A and determines the estimated time costs for sending and displaying the converted requests or the image at DE 278 A. The request times for the converted requests are added to the TRT value. Sample translations of X11 Windows display requests are shown in Table 1 of FIG. 5 along with a brief description of the operation. Given in Table 2 of FIG. 5 in columns 2 and 3 , are representative, empirically determined relative costs and estimated time costs (request times), respectively, for performing the operations in the converted request. At step 5 the converted display requests are placed at the head of Pending Request Queue (PRQ) 320 and process 310 repeats again from step 1 . At step 6 process 330 waits for the Total Network Time (TNT) for the selected us application to be less than a predetermined second value and for the total of the request times for the pending requests in PRQ 320 to exceed than a predetermined third value. If both conditions are met all pending requests in PRQ 320 are, at step 7 , optimized and written back into PRQ 320 and the TRT value and the TNT are adjusted based on the results of the optimization. The TNT represents the request times for all of the requests currently found in TRQ 340 . The optimization process of step 7 consists of step 7 A, prune optimization of the converted requests with those in PRQ 320 ; step 7 B merge optimization of the image with those requests in PRQ that effect the same screen area on DE 278 A; step 7 C, schedule optimization (attribute scheduling); and step 7 D, update optimization (update scheduling). The optimization process is further discussed below in the subsection entitled Queue Optimization. At step 8 the request at the tail of PRQ 320 is removed to be sent and a first sequence identifier, the send sequence number (SSN), is incremented and saved with the removed request to identify that request. At step 9 the removed request and its SSN are added to the head of Transmitted Request Queue (TRQ) 340 and the TNT is incremented by the request time of the newly added request (this also increments the TRT value). At step 10 the removed request is encoded into a form suitable for transmission over network connection P 7 and sent to DE 278 A. Preferably, to reduce the amount of data to be sent, the SSN identifier is not sent as part of the request. Requests are labeled or tagged with the SSN in the TRQ but the actual SSN identifier is not sent as part of the transmitted request. Corresponding identifiers are maintained by both the protocol engine and the display engine as each sends and receives requests. These identifiers are used as a reference point between the display engine and protocol engine. Again the order of steps 9 and 10 is not critical and the two operations can be performed in either order or simultaneously. The encoding process step 10 consisting of: step 10 A, coordinate compression; step 10 B, Run Length Encoding (RLE); step 10 C, serialization, step 10 D, delta compression; and step 10 E, Zlib compression. Encoding is further discussed in the subsection Encoding and Performance Monitoring set forth below. At step 11 the encoded request is written to the connection and sent via path P 7 to DE 278 A and process 330 repeats from step 6 . In DE 278 A, process 360 at step 12 is waiting for receipt of an encoded request. Upon receipt of an encoded request, at step 13 the request is unencoded consisting of one or more of the following depending on the compression techniques that were used in the request: step 13 A, Zlib expansion; step 13 B delta expansion; step 13 C, RLE expansion; and step 13 D, coordinate expansion. At step 14 the request is queued for execution and is displayed 372 on screen 370 of client device 214 . At step 15 a second sequence identifier, the Return Sequence Number (RSN), is generated by adding one to the current RSN and this will match the SSN identifier associated with the just executed request. At step 16 , if more than one second has past since the last time an RSN was sent or if there is no data now available from the network the current value of the RSN is sent to PE 276 A for use by process 350 . Process 360 then repeats beginning at step 12 . As is known, the protocol engine and display engine rely on the transport mechanism of the network connection P 7 to maintain the proper ordering or sequencing of the transmitted requests. This keeps the RSN identifier in synch with the SSN identifier and removes the need to send the SSN identifier with the transmitted request. At PE 276 A, process 350 at step 17 performs the RSN processing where the RSN identifier received from DE 278 A via connection P 7 is read and TRQ 340 is updated by removing each request from the tail of TRQ 340 until the SSN identifier for the request at the tail of TRQ 340 is greater than the RSN identifier and the TNT and TRT values are updated by subtracting the request time of each request as it is removed from TRQ 340 . These changes may result in the new total times falling below the predetermined first, second and third values mentioned above, in which case, more requests may be read and transmitted. At step 18 , process 350 performs periodic monitoring of the connection P 7 and adjusts encoding process 330 , as discussed below, in response to conditions (bandwidth available) found on the connection P 7 . Process 350 repeats starting at step 17 . Relative Cost and Request Times The cost of executing a display request on the client device consists of the following items: 1. Overhead of executing a single request irrespective of other factors. 2. Time cost based upon the number of pixels touched by the request. 3. An indication of the situations in which the request can be executed directly to the display hardware. If the request is not executable directly to the display then the display engine will simulate the request by drawing into memory and then copying this to the screen, thus incurring additional costs for the copying. The costs for 1 and 2 are a relative factor that is to be scaled by an overall performance figure calculated for the client device to provide a request time (see Estimated Request Time column in Table 2 of FIG. 5 ). As explained previously, the display engine executes a series of graphic requests both to the screen and in memory to compute this overall factor. This factor is also transmitted to PE 276 . Table 2 of FIG. 5 provides some samples of relative costs along with estimated request times for a specific client device. Table 2 of FIG. 5 provides relative costs—basic and per pixel—of the various graphics operations. Table 2 does not provide costs for all of the AIP request display operations. In actual operation, costs for all AIP requests supported by DE 278 are provided in a table in PE 276 . It contains a per pixel relative cost and a basic relative cost for each operation. The basic cost is empirically determined by testing execution of a given AIP request or display operation on a given class of client device in a given operating environment and is provided as part of the cost table that is found in each protocol engine. These basic costs can be thought of as being fixed. The performance factor—the per pixel cost, PPC—is determined as discussed previously and can be thought of as a dynamic or localized cost for the particular client device. The cost numbers are used by the cost calculation in the following way. The cost per pixel value is multiplied by an estimate of the number of pixels touched by the operation. The basic cost per display operation is multiplied by the number of basic operations that make up the request. For example, the operation AIP_POLYLINE is one where multiple line segments are being drawn so the basic relative cost for a 10 segment polyline operation would be 745 (10×74.50). Finally, the sum of these two values is multiplied by the performance factor (the PPC) that has been calculated for DE 276 on the client device to produce an approximation of the time taken by DE 276 to execute the request. The actual numbers in Table 2 of FIG. 5 are calculated from the following two equations: pixelCost=( m 2 −m 1 )( L 2 /b+L 1 /b )/ (n 2 −n 1 ) itemCost= m 1 −( L 1 /b )− n 1 *(sizeCost) These form the solution to the equations m 1 =L 1 /b+n 1 *pixelCost+itemCost m 2 =L 2 /b+n 2 *pixelCost+itemCost where m 1 and m 2 =single request time (s) for samples 1 and 2 L 1 and L 2 =single request length (bytes) for samples 1 and 2 b=client device bandwidth (bytes/s) N 1 and n 2 =number of pixels for samples 1 and 2 The results are computed using two samples per request tropically a 10 pixel test and a 500 pixel test. The results of the above calculations should then be multiplied by the current display rate of the test client device to produce the relative basic cost value shown in Table 2 of FIG. 5 . The estimated request times given in Table 2 of FIG. 5 are calculated for a client device that reports a performance figure of 9.15 megapixels per second. Bandwidth and Latency PE 276 A and DE 278 A cooperate to establish the bandwidth of connection P 7 and the latency. The latency is established as half of the average time that it takes to send a very small data packet from PE 276 A to the DE 278 A and back. The bandwidth is estimated by sending a series of large data packets and timing the difference between the arrival of the first packet and the last. The bandwidth is then defined as the total number of bytes transmitted divided by the total time difference. The bandwidth and latency are periodically measured. The size of the small request is about 3 bytes in length and the size of the large request is about 1024 bytes. The actual number of packets used is increased up to the point that more than 250 ms are needed to perform for the overall test. So the test is run starting using one packet, then if 250 ms have not been used, a series of 2 packets is then sent and again if 250 ms have not be used the test is repeated with a series of 3 packets. The number of packets is increased until about 250 ms are used for the test. The series is normally 5 packets but on lower bandwidth connections a smaller number is used to limit overall time taken. This request time for a display operation which is added to the TRT value is based upon the following: 1. The size of the request in bytes as it will be transmitted via the network connection. 2. The time cost for executing one of these requests on the DE. 3. The time cost for this request based upon the number pixels it will change. 4. The time cost for copying the results to the actual display if this request can not be executed directly to the client display. The size of the request is estimated by taking the size of the actual data points to be transmitted and multiplying these by the current compression ratio that is being achieved by the network encoding layer. This, in turn, is divided by the current network bandwidth to give an estimate of the total time required to transmit this request. The costs for 2, 3 and 4 are calculated by reference to table 279 obtained from DE 278 at startup. The total of these three values is then multiplied by the current overall performance factor of the DE to give an estimate of the time taken to execute the request by the DE. This is then combined with the estimated transmission time to give an overall request time for that request. The Total Request Time (TRT) is calculated for all queued requests (both those in the PRQ and those in the TRQ) for each application. New display requests from the application are not accepted by the protocol engine when the TRT exceeds the first preset value, typically 5 seconds. The PE 276 will not accept further display requests from the application when the request times of the requests in the PRQ and TRQ exceed this value so that the connection to the DE will not be overloaded. The Total Network Time (TNT) is comprised of the request times of the requests in TRQ 340 plus the latency time associated with the network connection. The PE maintains a TNT required to send all of the requests currently in TRQ 340 . The PE monitors the time associated with requests in the TRQ and will not transmit any further requests once this value reaches a pre-set value. This preset value is normally set at 2 seconds. This mechanism is used to ensure that a large backlog of requests does not build up for slowly executing client devices while ensuring that the full bandwidth of the network is used when this is the limiting factor. Queue Optimization Request optimization (FIG. 4, step 7 ) will take place if a reasonable number of requests are in the PRQ which is the case in normal operation. This process operates on an estimate of what will happen in the future in terms of request execution time not on what has happened in the past. This is important when dealing with requests for which the execution time varies greatly. When there is space to send new requests to the DE a process of optimizing the PRQ is performed. Basically this portion of the PE process looks through the requests in the PRQ. It attempts to remove requests that are redundant, converts sequences of requests that operate on the same screen area to a single less expensive request, re-orders the queue to reduce overhead in the DE and adds hints to the queue to make the performance of the DE more efficient. Because this process can be expensive in terms of time, the optimization is only performed when the total execution time of the requests in the PRQ exceeds a preset threshold that is proportional to the predetermined first value (typically 5 seconds). The typical value used for the predetermined third value is {fraction (1/20)} of the first value. These figures are all arbitrary and are based upon an acceptable level of application latency. Advantageously, for slow clients or networks more optimization is performed (since the cost of requests is higher) and that for faster networks and clients the optimization is only performed when there is a large amount of output. Preferably in step 7 of FIG. 4, the following four types of optimization are performed. These are pruning, merging, update scheduling (update optimization) and attribute scheduling (schedule optimization). The preferred method uses the following order of optimization—pruning, merging, following by, in either order, attribute scheduling and update scheduling. Pruning need not be done prior to merging but it is more efficient to do so. Update scheduling and attribute scheduling can be done in either order. Of the four types of optimization, merging is the most important in order to achieve good performance with the application and the connection. 1. In pruning (FIG. 4, step 7 A), converted requests that will be completely overwritten by a request later in the queue are removed. This optimization is equivalent to executing the series of graphics requests very quickly such that the effects of the earlier requests are not visible. Preferably, this is done prior to merging. 2. In merging (FIG. 4, step 7 B) if the cost of a group of converted requests that update a particular screen area is greater than the cost of transmitting an image of that portion of the screen then all of the requests are removed from the PRQ and an image request with its associated request time is inserted. The TRT is adjusted to reflect these removals and additions. 3. For update scheduling (FIG. 4, step 7 D), requests that can not be executed directly to the screen are grouped together. Hints are inserted into the request stream so that the DE can execute all of the requests in off-screen memory and then copy the results on screen in a single operation rather than requiring one copy operation per request. 4. For attribute scheduling (FIG. 4, step 7 C), requests which require the same color, and drawing modes are grouped together to avoid the need for the DE to switch back and forth for each request. This operation is equivalent to re-ordering the PRQ to make requests that share attributes adjacent within the PRQ. When performing this process it is necessary to ensure that the execution of overlapping requests is not changed. Encoding and Performance Monitoring The encoding process transforms a graphics request into a format suitable for transmission via the network connection to the DE. The DE inverts this process when the request is read from the network connection. The packet format used on the network connection takes the form of a header which has a byte indicating the type of the request, a two byte field which indicates the length of the packet and then a series of bytes which represent the actual request. The data for a particular request varies, but is, typically, either a series of coordinates for requests that involve some form of line drawing or a series of bytes for text and image requests. The encoding process is designed to enable the number of bytes required to represent the request to be minimized thus reducing the time required to transmit the data on a slow network connection. However, the processing time required to actually encode the data can be large. Thus, on higher bandwidth connections the AIP reduces the amount of encoding applied. This reduces the overhead but increases the number of bytes transmitted. The process of encoding uses the following steps 1. For coordinate compression, if the request contains a series of coordinates, all of these coordinates are converted to be relative to the first set in the request (FIG. 4, step 10 A). 2. If the request contains either text or an image it is compressed, preferably using the well know technique of Run Length Encoding (RLE) (FIG. 4, step 10 B). 3. The request is then serialized, i.e., converted into a simple series of bytes with a fixed byte ordering (FIG. 4, step 10 C). During this serialization process any small integer values in the range −64 to +128 are encoded into a single byte. 4. Requests that when encoded are less than 32 bytes in length are compared against a cache of previous requests. If the number of differences between the current request and a previous request is less than 8 bytes then the request is re-encoded as a reference to the previous request plus a series of two byte pairs describing the offset of the byte and the replacement value. This compression technique is known as delta compression (FIG. 4, step 10 D). If the encoded request is less than 32 bytes in length, it is entered into a delta compression cache (not shown) and the oldest item in the cache is discarded. Typically, the cache contains the last 32 requests although a lower or higher number of items can be chosen. 5. The encoded form of the request is compressed preferably using the well know Z-Lib compression algorithm (FIG. 4, step 10 E). 6. The encoded form of the request is sent over the connection to the DE (FIG. 4, step 11 ). At step 18 of FIG. 4, the AIP changes the encoding level based upon the bandwidth connection parameter of the connection as follows: 1. If the bandwidth is >300 Kbytes per second steps 3 and 4 are omitted. 2. If the bandwidth is <300 Kbytes per second steps 3 and 4 are applied. 3. If the bandwidth is >300 Kbytes per second the Z-Lib compressor is disabled. 4. If the bandwidth is between 100K and 300 Kbytes per second the Z-Lib compressor is set to use compression level 1 (the lowest and quickest to achieve compression level). 5. If the bandwidth is between 3 Kbytes and 100 Kbytes per second the Z-Lib compression level is set to level 5 (a mid range compression level). 6. If the bandwidth is less than 3 Kbytes per second the Z-Lib compression level is set to 9 (the highest compression level). Because the performance of the overall system can change over time AIP allows for both the available bandwidth of the network connection and the display speed of the DE to change. The PE uses two processes to track such changes. In the first performance updating process, the PE periodically re-checks the network bandwidth and latency. Typically, this is done about every five minutes. The PE also sends a request to the DE asking it to update the overall performance factor and send this new value back to the PE. These updated values are used in calculating the request times for newly added requests to PRQ 320 . Alternatively, if desired, the request times of requests already in PRQ 320 and TRQ 340 get updated. The details in how these values are obtained were discussed above in the description of the UAP system startup. In the second process, the PE when it moves a request from PRQ 320 to TRQ 340 adds to the request an estimate of the time the request should be removed from TRQ 340 . When the PE removes the request from TRQ 340 it maintains a running average and variance for the percentage error associated with this estimate. If the total of the average error minus half of the variance exceeds a pre-set limit (typically 200%) and the time since the last reevaluation is more than 1 minute then the performance updating process is invoked immediately. The overall effect of these processes is to provide a flow control mechanism that controls the rate that requests are read from the application and at which they are transmitted over the network to the client device. The goal of the AIP system is to ensure that DE 278 never stalls due to lack of available requests, that a large backlog of requests is not allowed to build up and that a reasonable number of requests are made visible to the UAP system within PRQ 320 since it is these requests that are available for optimization by the PE. It will be appreciated that for slower connections, requests will remain in the PRQ for a longer period of time so as to not slow down or overload the connection but advantageously giving the PE greater opportunity for optimization of the requests in the PRQ and conversely, when a more robust connection is available more requests can be sent utilizing the bandwidth of the connection while reducing the time spent in the PRQ increasing performance at the client device thus making the best possible use of the capabilities of the client device and the network connection for maintaining and enhancing performance. The implementation languages that are used are such that the protocol engines are typically created as a native binary for performance reasons. However, they could also be implemented in the Java programming language. The display engines are implemented in several languages to provide a broad coverage of client devices. The Java programming language is used to give support for platforms supporting either Java runtimes or Java supported in a web browser. However for systems that do not support a suitable Java environment native binary versions of the display engine components are used. Further, not all engines need to be in operation at all times, For example, if no client devices are requesting any application services, no protocol engines or display engines will be operating. Other embodiments of the invention will be apparent to those skilled in the art from a consideration of the specification or from practice of the invention disclosed herein. It is intended that the specification be considered as exemplary only with the scope and spirit of the present invention being indicated by the following claims.

In a three tier client-server architecture, an universal application (UAP) server providing a middle or second tier between a third tier of application servers and a first tier of client devices to provide the controlled delivery of computer based application programs running on a set of application servers to a number of users using a wide variety of client devices without requiring the installation of software on, or changes to the application servers or application programs. The UAP server or application broker comprises a status manager engine, a session manager engine, a data store, a data store engine, a web server, and multiple instantiations of protocol engines design for a given application type and multiple instantiations of display engines for downloading to and operating on client devices. The UAP server provides: (1) means to support standard graphics based computer applications connected to clients of varying capability via a network of varying bandwidth and latency by automatically varying the type and number of graphic requests and their networking encoding to provide near optimum performance while maintaining the correct visual representation, (2) automatic construction of user worktops based on user data contained in a datastore, (3) loading balancing of application servers upon application launch and (4) resumable connections to client devices.

REFERENCE TO RELATED APPLICATIONS [0001] This patent application claims the benefit of U.S. Provisional Application No. 61/318,781 filed on Mar. 29, 2010, the disclosure of which is incorporated herein by reference in its entirety. BACKGROUND OF THE INVENTION [0002] 1. Field of the Invention [0003] The present invention relates generally to an airship and more specifically to a system for actively controlling the lift of an airship. Lift is accomplished by a first gas such as air which is located in the airship at an internal pressure which is greater than the atmospheric pressure of the air on the outside of the airship and a bag filled with helium which is located within the airship and is surrounded by the first gas. The bag located within the airship and surrounded by the first gas has enough helium in it at ambient temperature to lift all but the cargo and fuel in the airship. When the first gas and helium are both heated by a heating means in the airship, increased buoyancy due to the increase of expansion of both heated gases will lift the airship with its fuel and cargo. [0004] 2. Description of Related Art [0005] Airships are known in the prior art. More specifically, by way of example, U.S. PreGrant Publication No. 2007/0102571 to Colting discloses an airship for lifting heavy and/or oversized loads. The airship uses the leverage of positive buoyancy to lift and transport payloads. [0006] U.S. PreGrant Publication No. 2005/0211845 to Perry; et al. discloses a non-rigid or semi-rigid airship with a hull having a plurality of lobes formed therein. The lobes decrease the radius of curvature of the hull, thereby reducing the stress on the hull due to the pressurized lifting gas contained therein. The reduced stress allows the hull to be constructed from a lighter weight material, thus reducing the mass of the hull, and enabling the airship to carry more cargo. Flexible membranes are used to partially delineate lobes. The membranes are attached to the inner surface of the hull and a group of load lines connected to and running between the membranes form a polygon-shaped cross-sectional area. [0007] U.S. Pat. No. 7,866,601 to Balaskovic discloses an airship shaped as an oblate spheroid and a support structure which forms a partial support for the hull. A horizontal stabilizing member is coupled to a lower surface of the airship, and a vertical stabilizing member having a first end is pivotally coupled to the airship. The vertical stabilizing member and the horizontal stabilizing member may be operably coupled to one another. [0008] U.S. Pat. No. 7,156,342 to Heaven, Jr., et al. discloses a system for actively controlling the aerostatic lift of an airship by manipulating the ratio of air to lifting gas contained within the airship, and thus the overall mass of the airship. This manipulation is accomplished by actively compressing and/or decompressing the lifting gas or internal air, with the resulting pressure differential borne primarily by the hull and/or an internal pressure tank. [0009] U.S. Pat. No. 6,837,458 to Swearingen, et al. discloses an airship having a hull which includes a first section having a width which varies along the selected direction of travel where the width increases from the bow of the hull to a maximum width and then decreases from the maximum width to the tail section of the first section; and a second section coupled to the first section and having a width which varies along the selected direction of travel where the width increases from a leading edge of the second section to a maximum width and decreases from the maximum width to the stern of the hull. [0010] U.S. Pat. No. 6,793,180 to Nachbar, et al. discloses an airship hull having a plurality of flexible members disposed lengthwise about the perimeter of the airship skin. The flexible members can be held in place in sleeves on the skin of the airship where there ends are drawn toward one another by tensioning means which cause the members to bow outwardly from a central axis to provide a rigid structure for the skin. [0011] U.S. Pat. No. 6,293,493 to Eichstedt, et al. discloses a non-ridged semi-buoyant vehicle with a pressure stabilized gasbag which has an aerodynamic shape. The gasbag includes vertical catenary curtains, a pair of first and second Y shaped catenary curtains which are coupled to the vertical catenary curtains and extend along a second portion of the gasbag and the arms of each of the Y shaped curtains are attached to the top surface and the legs are attached to the bottom surface of the gasbag. [0012] U.S. Pat. No. 5,890,676 to Coleman, et al. discloses a neutral buoyancy fuel bladder which uses hydrogen and oxygen to power an airship. The neutral buoyancy fuel bladder includes a fuel cell, electrolyzer, and means for storing hydrogen, oxygen and water. The fuel cell uses the hydrogen and oxygen to create heat, water and current flow. An energy source transmits a beam to an energy receiving unit on the airship, and the energy from said beam is used to power said airship, and replenish the supply of hydrogen and oxygen. [0013] U.S. Pat. No. 4,591,112 to Piasecki, et al. discloses an airship with provisions for vectored thrust provided by a plurality of controllable pitch rotor thrust producing units attached to the hull. The control systems are interconnected to be operable by a master control which establishes both similar and differential pitch settings of the rotors of selected thrust units in a manner to establish vectored thrust in directions which establish the required amounts of vertical lift, propulsion thrust, trim and control forces to control all flight aspects of the airship. [0014] U.S. Pat. No. 4,326,681 to Eshoo discloses a lighter-than-air disc-shaped non-rigid airship having a flexible envelope within which an annular pressurized tube is positioned to maintain the flexible envelope in a saucer shape when inflated. A gondola is suspended beneath the central chamber. To maintain level horizontal flight stability, differential forces are developed by providing the central chamber with heated air and the outer chambers with a lighter-than-air gas such as helium to provide greater lift than the central chamber. Propulsion units are arranged at opposite side edges of the envelope and maneuvering is accomplished by rotating the airship. [0015] Air vehicles that use gas that is lighter than air as a form of buoyancy control have been successfully flown for over 200 years. Common gasses utilized are helium, hot air and hydrogen. In recent years there have been many attempts to design a practical cargo carrying air vehicle that uses buoyant gas to aid in lifting the vehicle and cargo. There has been little or no success in these attempts. [0016] The main feature required for a successful cargo carrying air vehicle that uses buoyant gas is that it needs to have the ability to vary its buoyancy in order to pick up and drop off cargo and to reduce its buoyancy as fuel is burned off. [0017] Most applications for buoyant cargo carrying air vehicles require that the cargo be dropped off at its destination and that the vehicle then return home without any cargo. This requires that the air vehicle have the ability to vary its buoyancy by at least the amount of the cargo weight. [0018] It has proven impractical in most cases to replace the cargo with disposable ballast like water for the return trip. Large quantities of water are not always available at the destination. [0019] There is a solution for the fuel burn-off buoyancy problem but it involves using complicated exhaust condensation devices or nonstandard fuels such as blau gas. [0020] Recent attempts to solve this buoyancy problem have included using helicopter like rotor lift to carry the payload and fuel, using aerodynamic lift that is generated by forward motion and helium recompression. [0021] Rotor lift buoyancy control consumes fuel at a very high rate and is not practical for long distance use. [0022] Using aerodynamic lift from forward motion negates one of the main advantages of buoyant gas vehicles which is the ability to take off and land vertically. Also, the aspect ratio of any wing type surface will be so low that this will again require tremendous amounts of fuel to achieve the necessary lift. [0023] Helium recompression equipment is very heavy and the process is too slow to be practical. [0024] Pressurized and powered hot air vehicles, like hot air balloons, have recently been developed. Size and significant fuel burn has limited their success. Very large envelopes are required as hot air lifts only about 20 to 25% of what the same volume of helium lifts. Also, these large envelopes radiate a lot of heat so that fuel burn is great. [0025] My air vehicle design overcomes the buoyancy problem completely and efficiently. [0026] The other problem area for buoyant air vehicles is dealing with size. Any air vehicle that uses a buoyant gas to aid in lifting payload must, by their very nature, be very large. Large vehicles are difficult to deal with when a wind storm arrives. [0027] If hot air is the buoyancy source, you can release the hot air, fold up the envelope and seek shelter from the storm. This is not practical for a large commercial cargo carrying vehicle. [0028] If helium or hydrogen is the lifting gas, it is too expensive and impractical to either re-pressurize the gas into high pressure containers or to vent the gas off into the atmosphere. Even if the gas were removed, there would still be a lot of envelope lying on the ground that must be secured. [0029] The only solutions found to date to secure an airship during inclement weather are to store it in a hangar or to secure the nose of the airship to a mooring mast in an area large enough to let the airship weathervane in all directions [0030] Large airship hangars are very expensive and prove impractical for that reason. A mooring mast is also expensive and they must be large, permanent structures for large airships. [0031] Mooring masts do not provide the required level of protection for an airship. Many airships have been destroyed while on mooring masts in less than extreme weather. Mooring mast damage can come from many sources. Gondola damage occurs when vertical wind gusts raise and lower the airship. There have been cases where airships were raised vertically above the mast before they were forcibly returned to the ground. Airship envelopes have been torn apart by the stresses on their single point nose attachment to the mooring mast. [0032] My air vehicle design overcomes the mooring problem and can be moored from a single point on the ground in winds exceeding 100 mph. [0033] Another problem area for non-rigid airships is the need to make the nose of the airship less prone to implosion as a result of dynamic air pressure when moving at high airspeeds. The problem is generally addressed with nose battens which are heavy and difficult to deal with. [0034] My air vehicle design addresses this problem in a way that will allow my design to fly faster than normal non-rigid airships. SUMMARY OF THE INVENTION [0035] My design for a cargo carrying air vehicle includes an outer insulated pressurized air envelope made out of flexible material with a limp internal helium filled bag. This internal bag has enough helium in it at ambient temperature to lift all of the structure but not the fuel or cargo. This is the reverse of a normal airship design where the outer envelope holds the helium and then the pressurized helium envelope contains air ballonets inside it. [0036] In my design, the limp helium bag is large enough to allow the helium to expand to 1.5 to 3 or more times its ambient volume before it reaches a pressurized state where helium must be vented from it to preserve its integrity. This number depends on the size and purpose of the vehicle being designed. The helium and/or air inside this outer insulated envelope is directly heated with a furnace to a maximum operating temperature, in the neighborhood of 300 degrees Fahrenheit. This heated air and helium will pick up the fuel and cargo due to the increased buoyancy of both hot expanded gasses. Then with normal airship thrusters the air vehicle will transport the cargo at a very low cost. Because of the ratio of allowed helium expansion before venting, this design can attain altitudes of over 25,000 feet which was not achievable with cargo carrying airships of the past. [0037] The main envelope is pressurized with air to a normal non-rigid airship pressure. The shape of this envelope is a generally typical streamlined airship shape. This envelope has to be designed to prevent air and heat from leaking through it. The easiest way to prevent heat from leaking out will be as simple as affixing ordinary un-faced fiberglass house insulation or other material to the inside of the main envelope. [0038] Depending on specific size and design of the vehicle, the helium quantity in the helium bag may be from about 20% to 40% of the volume of the main envelope. In normal operation this helium bag will never contain any significant pressure and this reduces its weight and cost. [0039] The helium membrane needs only to be capable of holding helium, to transfer heat efficiently and withstand temperatures of approximately 300 degrees Fahrenheit. These requirements are achieved with the same material, a vacuum deposited aluminum on a flexible and thin high temperature substrate. One likely substrate material is a high melt temperature nylon. [0040] The envelope furnace heats the inner air and/or the helium directly to a normal operating temperature of approximately 250 to 300 degrees Fahrenheit or higher. [0041] A gondola and/or separate cargo carrier is hung below the main envelope on long cables/ropes. These long cables are necessary to maintain pendulum stability of the air vehicle. The air vehicle is designed to always fly at a level attitude. Going up or down is achieved by add or removing heat from the envelope, not by raising or lowering the nose as an aircraft or blimp does. [0042] These long cables also provide the solution to the airship mooring problem. The gondola can be moored tightly to the ground by a single point and the envelope will weathervane above it at a high altitude. Since the helium is always lifting at least the weight of the whole structure, the gondola cables will always have as much tension as the weight of the gondola. This tension will keep the weather vaning envelope high overhead in winds up to about 70 or 80 mph. In winds higher than about 80 mph, heat is added to the envelope to put more tension on the gondola cables. This forces the envelope to ride higher in the air to a more stable position. Winds above 70 or 80 mph are generally never sustained very long and therefore minimal fuel would be required to stabilize the vehicle until the storm subsides. [0043] The method that I use to prevent envelope nose implosion at high speeds is to use what I refer to as ball-cones. These pressurized, air filled nose envelopes whose air pressure is many times that of the main envelope, form a more pointed and much stiffer nose and tail structure for the air vehicle. [0044] Propulsion of the vehicle will be provided by thrusters on the envelope and/or gondola/cargo carrier. Yaw control of the vehicle can be provided by movable rudder surfaces or by angling the thrust line of the propulsion motors. Precise control for hovering situations can be by any number of reversible and/or gimbaled thrusters attached to the vehicle. [0045] Vertical fins may be necessary to provide stability about the yaw axis. [0046] Reserve fuel may be kept on board for the amount of reserve time desired. Reserve fuel will be lifted by the ambient helium volume. This means that if it gets used on a flight, helium may need to be vented in order to land. [0047] My design becomes competitive with C130 cargo aircraft in depreciated dollars per ton per mile costs at approximately 380 feet in length. [0048] My 380 foot design will carry about 9,000 lbs of cargo 500 miles on about 470 gallons of fuel. [0049] My design's combination of using hot helium and hot air with virtually unrestricted helium expansion is actually a very efficient way to lift cargo weight. A normal cargo aircraft uses 30% to 40% of its fuel to lift its structure, fuel and cargo. My high speed large designs use as little as 5% of their total fuel to carry the structure, fuel and cargo. My slower smaller designs use up to 30%. [0050] My 800 foot design, the size of the Hindenburg, will compete directly with the largest transport aircraft, the C5 Galaxy, in terms of depreciated dollars per ton per mile costs. [0051] My 1100 foot to 1500 designs will satisfy the proposed 21 st century Global Range Aircraft request of the US military. That is that they can deliver 150 tons of cargo 12,000 nautical miles and then return home without cargo and without refueling. [0052] To be clear on the operation of my air vehicle, this is a typical flight profile of Version 2 of my design as depicted in FIG. 31 [0053] The pilot arrives at the vehicle mooring location and begins her routine. As she approaches the vehicle she hears the sound of one of the on board generators running. This is good because she knows it is likely that the cargo carrier's air filled structure and vehicle's 1,075 foot long envelope, way up in the sky, will likely be at their operating air pressure which must be maintained 24 hours a day. [0054] When she arrives in the cockpit, she confirms that all the pressures are OK. Next, she depressurizes the cargo carrier, opens the access doors and lowers the loading ramps. [0055] At this point the cargo handlers busily begin loading the cargo carrier with today's cargo while the pilot finishes her daily inspections. She then calculates the amount of fuel that needs to be carried in the thirty one, 1,000 gallon tanks that are in the floor of the cargo carrier's deck. Since the cargo weight and the fuel weight need to be lifted into the air with heat energy, there is no use carrying more fuel than necessary for the today's trip. [0056] The load today is 275 tons of mining equipment that needs to be delivered to a diamond mine that is 250 miles away. The vehicle will then return home, empty of cargo, without refueling. At an 80 mph cruise speed, that will require about 3400 gallons of fuel. No need to worry about adding any reserve fuel as 500 gallons of reserve fuel is always carried aboard on the engine gondola hanging from the envelope cables 100 feet above her cockpit. The reserve fuel should never be touched in normal operation so its weight has to be added to the weight that the ambient temperature helium has to lift. This means that if the reserve fuel is ever used, the pilot would have to vent off helium to land, if the on board cargo weight at the time, is less than the weight of the reserve fuel used to get to the destination. [0057] About an hour before the loading is finished, the pilot turns on the envelope's furnace to bring the internal envelope temperature to 300 degrees. Since the vehicle is lifting a large load today, it will take about an hour for the 100 million BTU furnace to heat the envelope gasses up to a temperature that will lift the vehicle with its cargo and fuel into the air. Since the vehicle is held to the ground by its single mooring cable, there is no worry that the vehicle will lift off prematurely. [0058] With the cargo loaded and the envelope temperature is within 5 minutes of reaching its 300 degree liftoff temperature the pilot closes the ramps and access doors. Then she pressurizes the cargo carrier with air to make it rigid enough to handle the winds at cruise speed. Next, the pilot starts the four 2,000 horsepower propulsion motors. [0059] By the time the engines are running, the mooring cable begins to tighten and the pilot casts off the mooring line. The vehicle begins rising into the air. She sets all four engines to cruise power and then turns the vehicle to the cruise heading with her rudder pedals. The rudder pedals command the fore and aft envelope engines to swing to the left which turns the vehicle. [0060] At this point the pilot increases the commanded envelope temperature to 325 degrees, its maximum operating temperature. This causes the climb rate to increase to 500 feet per minute. At about 2,000 feet the climb rate goes down as the vehicle is reaching the maximum height that 325 degrees will elevate it to today with this large cargo weight and fuel. With less payload it could climb much higher. The vehicle ceiling with no cargo and low fuel is about 28,000 feet. [0061] As the trip progresses, fuel is burned off which allows the vehicle to gradually climb higher. In about two and a half hours it has reached 3,500 feet but it is time to descend. The pilot now switches the furnace control to Climb/Descent Mode and commands a 100 foot per minute rate of descent which should get the vehicle to the mine site at about 500 feet above the drop off point. With this 100 foot per minute descent command, the onboard computer modulates the burner flame intensity to achieve the commanded rate of descent. Although the pilot could easily do this herself, this is a perfect job for a computer and the pilot can do other more important things while descending. [0062] Near the drop off point the pilot switches the burner control to manual. This locks the burner flame size at its current setting which was achieving the 100 foot per minute descent. The final landing position control will use gimbaled engine thrust. [0063] The pilot then lowers the mooring cable to hang 200 feet below the cargo carrier and hovers the cargo carrier near the ground mooring point using the gimbaled fore and aft engines. The joystick in the pilots hand controls the gimbaled engines. The vehicle responds by moving left/right for left/right joystick movement, forward/backward for forward/backward joystick movement, up/down for up/down force on the joystick and a left/right yaw for a left/right twisting force on the joystick. The more joystick displacement left/right and forward/backward or force that is applied to the joystick for yaw and up/down, the more pitch that is applied to the appropriate thrusters. When the pitch reaches the maximum for that propeller, the power level of that thruster is increased up to its maximum power level. Each joystick direction has a trim slide that varies the appropriate thruster forces when the pilot is not holding the joystick. [0064] With this precise control available a ground worker can safely pickup and attach the end of the mooring cable that was extended from the cargo carrier to the ground mooring point. The pilot then winches the vehicle to the ground in the exact location of the ground mooring point. A few minutes before this, the pilot depressurized the cargo carrier so that the loading ramps could be lowered immediately after landing and the access doors opened. Unloading now begins. [0065] As soon as the cable was attached to the ground and the slack removed, the pilot turned off the burner. This began the cool-down that is needed before the vehicle can be released from the ground after the cargo has been unloaded. [0066] Because the pilot knows this ground crew will unload the vehicle quickly, the pilot increases the cool-down rate by opening the aft cool down vents on the envelope and places the aft engine in full reverse thrust. This blows cold air through the envelope, forcing hot air to exit the front overpressure vents. This forced cooling of the envelope air can reduce the buoyancy of the vehicle faster than the weight can be removed from the cargo carrier by the ground crew, in fact, in less than half an hour for this model of air vehicle. In an hour the cargo is unloaded and the vehicle is ready to takeoff for home. The pilot now closes and pressurizes the cargo carrier. [0067] For the return flight the pilot decides to use Climb/Descent Mode on the furnace control since the vehicle is light now and the cruising altitude will be higher. To liftoff with no cargo and only the remaining fuel, the envelope temperature needs only to be a few degrees above ambient, so very little burner will be used. This means the vehicle will be flying high and the pilot will need the pressurized cockpit for the flight home. [0068] To takeoff, the pilot starts the four thruster motors, releases the mooring line from the ground mooring and sets the commanded climb rate to 2000 feet per minute. The vehicle rises within a minute and the pilot heads for home. After 25 miles the vehicle is at 27,000 feet and cruising at about 110 mph. This is the ceiling or pressure altitude with this amount of fuel and cargo. This means that the vehicle is cruising with the helium expanded to near its maximum volume and the burner control has sensed this. The burner control is now modulating the flame based on keeping the helium volume at its maximum operating limit. [0069] About 50 miles from home the pilot selects a 500 foot per minute descent and by the time she is home the vehicle is a couple hundred feet high and ready to land. But, there is no ground person around that can grab the mooring line and attach it to the ground mooring point. This is no problem, the pilot simply lands near her ground mooring point with her joystick thrust control and makes sure that the envelope is cool enough to stay put on the ground for a minute while she gets out and attaches the vehicles mooring cable to the ground mooring point. She then gets back in the cockpit and uses the joystick thrust or a little burner burst to get back in the air while the vehicle winches itself to the ground. [0070] As the pilot shuts down, she notes that there is about 400 gallons of fuel left. She also notes that the whole trip took about 8 hours. Before she leaves, she makes sure that the generator is running and that all pressures are OK. [0071] In an exemplary embodiment of the present invention, there is disclosed system for controlling the lift of an airship for carrying a cargo and a supply of fuel comprising: [0072] a self supporting hull made of a flexible gas impermeable material; [0073] a first gas located in the hull at an internal pressure which is greater than the atmospheric pressure of the air on the outside of the hull; [0074] a bag filled with helium located within the hull and surrounded by the first gas; [0075] a heating means coupled to heat the first gas and the helium; [0076] at least two airship thrusters coupled to urge the airship to travel to a selected destination using the fuel carried by the airship; and [0077] a gondola coupled to the self supporting hull; [0078] wherein the bag located within the hull and surrounded by the first gas has enough helium in it at ambient temperature to lift all but the cargo and fuel in the airship; [0079] wherein the first gas and helium when heated provides increased buoyancy due to increased of both expansion of both heated gases to lift the airship with its fuel and cargo. [0080] The more important features of the invention have thus been outlined in order that the more detailed description that follows may be better understood and in order that the present contribution to the art may better be appreciated. Additional features of the invention will be described hereinafter and will form the subject matter of the claims that follow. [0081] Before explaining at least one embodiment of the invention in detail, it is to be understood that the invention is not limited in its application to the details of construction and the arrangements of the components set forth in the following description or illustrated in the drawings. The invention is capable of other embodiments and of being practiced and carried out in various ways. Also it is to be understood that the phraseology and terminology employed herein are for the purpose of description and should not be regarded as limiting. [0082] As such, those skilled in the art will appreciate that the conception, upon which this disclosure is based, may readily be utilized as a basis for the designing of other structures, methods and systems for carrying out the several purposes of the present invention. It is important, therefore, that the claims be regarded as including such equivalent constructions insofar as they do not depart from the spirit and scope of the present invention. [0083] The foregoing has outlined, rather broadly, the preferred feature of the present invention so that those skilled in the art may better understand the detailed description of the invention that follows. Additional features of the invention will be described hereinafter that form the subject of the claims of the invention. Those skilled in the art should appreciate that they can readily use the disclosed conception and specific embodiment as a basis for designing or modifying other structures for carrying out the same purposes of the present invention and that such other structures do not depart from the spirit and scope of the invention in its broadest form. BRIEF DESCRIPTION OF THE DRAWINGS [0084] Other aspects, features, and advantages of the present invention will become more fully apparent from the following detailed description, the appended claim, and the accompanying drawings in which similar elements are given similar reference numerals. [0085] FIG. 1 shows a perspective view of Version 1 of the air vehicle with a cargo load slung below it. [0086] FIG. 2 shows a top view of Version 1 of the air vehicle. [0087] FIG. 3 shows a side view of Version 1 of the air vehicle. [0088] FIG. 4 shows a front view of Version 1 of the air vehicle. [0089] FIG. 5 shows Version 1 of the air vehicle moored to the ground in light winds. [0090] FIG. 6 shows Version 1 of the air vehicle moored to the ground in medium winds. [0091] FIG. 7 shows Version 1 of the air vehicle moored to the ground in high winds. [0092] FIG. 8 shows Version 1 of the air vehicle moored to the ground in preparation for winching the envelope to the ground mooring position. [0093] FIG. 9 shows Version 1 of the air vehicle with the gondola removed and the envelope winched halfway to the ground mooring position. [0094] FIG. 10 shows Version 1 of the air vehicle moored to the ground in the ground mooring position. [0095] FIG. 11 shows an exploded x-ray view of envelope components of Version 1 of the air vehicle. [0096] FIG. 12 shows an x-ray view of the envelope components of Version 1 of the air vehicle. [0097] FIG. 13 shows an exploded x-ray view of separate gas areas of Version 1 of the air vehicle. [0098] FIG. 14 shows an extracted view of the contracted helium bag of Version 1 of the air vehicle at ambient temperature. [0099] FIG. 15 shows an extracted view of the fully expanded helium bag of Version 1 of the air vehicle at an elevated temperature and/or reduced pressure. [0100] FIG. 16 shows a flat view of the helium bag of FIGS. 14 and 15 compared to a side view of the main envelope. [0101] FIG. 17 shows a view of the contracted helium bag of Version 2 of the air vehicle at ambient temperature with the main envelope cut away. [0102] FIG. 18 shows a view of the fully expanded helium bag of Version 2 of the air vehicle at an elevated temperature and/or reduced pressure with the main envelope cut away. [0103] FIG. 19 shows an x-ray side view of the envelope structure of Version 2 of the air vehicle with the helium bag contracted. [0104] FIG. 20 shows an x-ray side view of the envelope structure of Version 2 of the air vehicle with the helium bag fully expanded. [0105] FIG. 21 shows a cut away view of the furnace for Version 2 of the air vehicle. [0106] FIG. 22 shows an exploded view of the burner section of the furnace for Version 2 of the air vehicle. [0107] FIG. 23 shows an exploded view of a section of the heat exchanger in the air portion of the furnace for Version 2 of the air vehicle. [0108] FIG. 24 shows an exploded view of a section of the heat exchanger in the helium portion of the furnace for Version 2 of the air vehicle. [0109] FIG. 25 shows an exploded view of the helium fan section of the heat exchanger in the helium portion of the furnace for Version 2 of the air vehicle. [0110] FIG. 26 shows an x-ray view of the furnace for Version 2 of the air vehicle. [0111] FIG. 27 shows an exploded view of the separator section of the furnace for Version 2 of the air vehicle. [0112] FIG. 28 shows the envelope cooling entrance ducts in an open position on an isolated section of the rear of main envelope. [0113] FIG. 29 shows the envelope cooling entrance ducts in a closed position on an isolated section of the rear of main envelope. [0114] FIG. 30 shows the envelope cooling entrance ducts in an open position on the rear of main envelope and the cooling exit ducts in an open position on the front of the envelope. [0115] FIG. 31 shows Version 2 of the air vehicle with 4 engines and a cargo carrier structure slung below the gondola. [0116] FIG. 32 shows the gondola and cargo carrier of Version 2 of the air vehicle, moored on the ground, with the loading ramps down and the front and rear access doors open. [0117] FIG. 33 shows a bottom perspective view of the gondola and cargo carrier in the air. [0118] FIG. 34 shows a bottom perspective view of the gondola and cargo carrier in the air with the landing cushion membrane removed. [0119] FIG. 35 shows a bottom perspective view of a single pallet from the pallet array that composes the load carrying section of the cargo carrier of Version 2 of the air vehicle. [0120] FIG. 36 shows a cut away bottom perspective view of the pallet described in FIG. 32 . [0121] FIG. 37 shows an exploded lower view of the pallet described in FIG. 32 . [0122] FIG. 38 shows a perspective section view of three side trusses in the pallet array described on FIG. 32 . [0123] FIG. 39 shows an isolated section view of the spring end of the joining cables for the pallet array described in FIG. 32 . [0124] FIG. 40 shows an isolated section view of the swaged cap end of the joining cables for the pallet array described in FIG. 32 . [0125] FIG. 41 shows Version 2 of the air vehicle on a single point ground mooring, in high winds that have shifted 45 degrees since the vehicle landed. [0126] FIG. 42 shows an expandable streamlining cuff assembly for the gondola and cargo carrier support cables. [0127] FIG. 43 shows the internal load transfer structure of the gondola. [0128] FIG. 44 shows a front skeleton view of the cargo carrier with the fabric cover removed, the loading ramps down and the front access door open. [0129] FIG. 45 shows a rear skeleton view of the cargo carrier with the fabric cover and hoops removed, the loading ramps down and both front access doors closed. [0130] FIG. 46 shows Version 2 of the air vehicle having just landed at a site with nine mooring points that allow full maintenance of the vehicle. [0131] FIG. 47 shows Version 2 of the air vehicle repositioning itself to the main single point mooring location of the site described in FIG. 46 . [0132] FIG. 48 shows Version 2 of the air vehicle with its envelope maintenance mooring lines lowered and attached to their ground winches at the site described in FIG. 46 . [0133] FIG. 49 shows how the gondola is lowered into a maintenance saddle as the envelope is winched to the ground at the site described in FIG. 46 . [0134] FIG. 50 shows Version 2 of the air vehicle fully secured in its maintenance position with the lower portion of the envelope squashed securely onto the ground. [0135] FIG. 51 shows a joystick control to operate thrusters DESCRIPTION OF THE PREFERRED EMBODIMENT [0136] FIG. 1 depicts Version 1 of the air vehicle with a cargo load slung below it. Version 1 is a 240 foot long envelope 1 whose displacement is 250,000 cubic feet. To provide positive yaw stability, two vertical surfaces, 2 and 3 have been attached to the rear of envelope 1 . A propulsion unit 4 has been attached to the rear on envelope 1 . Propulsion unit 4 is mounted on a swivel to allow it to be angled to the left and right of the vehicle. In doing this, propulsion unit 4 provides yaw control of the vehicle in addition to its use for forward thrust. It is likely that all propulsion units have the ability to provide thrust in both directions through their propeller. Hanging from the envelope 1 are four primary gondola cables 6 . The forward two primary gondola cables 6 are joined together at their ends. From this joining point, a secondary gondola cable 7 continues down to the gondola 8 . The two rear primary gondola cables 6 are similarly joined and attached to another secondary gondola cable 7 . In this Version, gondola 8 is simply a Cessna 172 aircraft fuselage with the wings and horizontal tail surfaces removed. This fuselage, with its motor and landing gear provides all of features needed in our gondola 8 and is a relatively inexpensive and quick solution for a gondola on this size air vehicle. [0137] Slung below the gondola 8 is a cargo load 9 , which this airship can deliver. In this case the cargo is a Prius automobile. [0138] FIG. 2 shows a top view of Version 1 of the air vehicle. The main feature visible in this top view are two load bridles 5 that distribute primary gondola cable 6 loads over a large area of the envelope 1 . There are many other options for distributing the primary gondola cable 6 loads to the envelope 1 including load curtains or load patches. [0139] FIG. 3 shows a side view of Version 1 of the air vehicle. This provides a good depiction of how the two load bridles 5 distribute primary gondola cable 6 loads out over the top of envelope 1 in a fan shaped manner. [0140] FIG. 4 shows a front view of Version 1 of the air vehicle. This view again shows the load distribution of the primary gondola cables 6 using the load bridle 5 to transfer the loads from one primary gondola cable 6 , over the top of the envelope 1 and back to another primary gondola cable 6 on the other side. From this point on, load bridles 5 will be omitted from all drawings for the sake of clarity. [0141] FIG. 5 to FIG. 7 depict Version 1 of the air vehicle moored to the ground, in a field 11 , at a single point 12 in increasing winds. In the light winds of FIG. 5 , the envelope remains above the gondola 8 without any noticeable affects. [0142] In the medium winds of FIG. 6 , the envelope 1 will tend to tilt back a little. Because of the envelope's positive yaw stability, the envelope will weathervane into the wind. This tilted back angle presents no immediate stability problems until the angle increases to an unstable amount. [0143] In the strong winds if FIG. 7 , the front secondary gondola cable has been winched in by a winch mounted in the gondola 8 . This lowers the nose of the envelope 1 back to a horizontal and stable position. This winching in of the front secondary gondola cable will be activated automatically. All pressurized non-rigid airships, as this one is, require automatic 24 hour a day systems to ensure that the envelope stays pressurized. The power to maintain this comes from generators or plug in electrical cables. Since the power is there for a 24 hour a day system, this winch system for maintaining attitude in high winds in not a significant problem. As a last resort to keep the envelope 1 stable in high winds, the envelope heat can be automatically turned on. This will increase the buoyancy of the envelope 1 and the envelope will rise higher into the wind to a more stable attitude. [0144] If the wind changes direction a large amount, the secondary gondola cables 7 will begin twist as the gondola 8 remains pointed in it original direction. Larger versions of the design have a swivel to take care of this but on this size the solution is to simply add some heat to the envelope 1 which will lift the gondola 8 and then the gondola 8 will automatically spin around to untwist the secondary gondola cables 7 . At this point, the gondola 8 can be secured back down to the ground. [0145] FIG. 8 to FIG. 10 depict the procedure that will be followed to put the envelope 1 into a moored maintenance position with the envelope 1 securely fastened to the ground. In this state, the envelope 1 mimics an inflatable ground structure. [0146] In FIG. 8 , the gondola 8 has landed in the middle of four ground anchor point/winches 17 . Here, the maintenance winch lines 13 and 14 are introduced. These lines are lowered from winches attached to the rear and lower middle of envelope 1 . They are controlled by remote control or hard wired lines that run up the secondary and primary gondola cables 6 and 7 . The purpose of the maintenance winch lines 13 and 14 is to allow a person to be raised up to the envelope 1 in a basket 15 to do maintenance on the propulsion unit 4 and FIG. 17 furnace 23 . Maintenance winch lines 13 and 14 are also used to lower the envelope maintenance cables 16 from their stowed position. Once the maintenance cables 16 are lowered and attached to the four ground anchor point/winches 17 , the envelope 1 can be winched down to a point where the gondola 8 can be disconnected from the secondary gondola cables 7 and then the gondola 8 can be moved out from under the envelope 1 . FIG. 9 shows the envelope 1 in this position with the gondola 8 removed. [0147] FIG. 10 shows the envelope 1 fully winched down and secured in its maintenance mooring position. Approximately 3,000 square feet of envelope 1 has been flattened out onto the ground. With this area of ground contact and with envelope 1 maintained at its normal flying pressure, the four maintenance cables 16 are under significant load and envelope 1 can withstand high winds from any angle. [0148] FIG. 11 to FIG. 13 depict the internal structure and components of envelope 1 . The basic elements of the flexible fabric elements of envelope 1 are shown in an exploded x-ray view in FIG. 11 . The components are the insulated outer fabric shell 18 , the helium bag 20 and the ball-cones 19 . [0149] The ball-cones 19 are a very significant design element. The shape of ball-cones 19 are spherical and slightly bigger than a half sphere on the end attached to the outer fabric shell 18 , and they have a basically conical shape that follows the contours of the FIG. 1 envelope 1 on the end away from the outer fabric shell 18 . At some point near the tip of the conical section, another, higher pressure ball-cone 19 can be added. In the drawings shown, the end of the conical sections are just a small half sphere closure and only one ball cone per FIG. 1 envelope 1 end is used. This is a way to create a very streamlined, rigid nose and tail section for air vehicle. The air pressure in the ball cones 19 shown will be at least four times the main FIG. 1 envelope 1 pressure. [0150] In FIG. 11 to FIG. 14 , the helium bag 20 is shown with about 100,000 cubic feet of helium in it. With this much helium in the bag and with the bag placed inside the pressurized and insulated outer fabric shell 18 , as in FIG. 12 , the helium bag takes on the shape shown. There is a lot of helium bag material unexpanded in the lower portion of the helium bag. [0151] FIG. 13 shows an exploded x-ray view of separate gas areas of Version 1 of the air vehicle. The volume of the heated air section 21 of the air vehicle is depicted at about 150,000 cubic feet. [0152] FIG. 15 shows the helium bag 20 expanded to its full size of about 165,000 cubic feet. This would be the condition normally seen at the pressure altitude of an airship. [0153] FIG. 16 shows the shape of the flat and empty helium bag 20 compared to the side view of an inflated envelope 1 . Because the helium bag 20 is just a flat shape, it will be easy to manufacture. This shape has also been tested for helium shifting when the pitch attitude of envelope 1 changes. The tests show that with the secondary gondola FIG. 8 cable 7 lengths shown, the airship will be stable in pitch without the need to restrict or secure the helium bag in any way inside the envelope 1 . [0154] Version 1 of the air vehicle will have the furnace heat only the heated air FIG. 13 section 21 of the envelope 1 . The helium with be heated and cooled by the heat transfer across the helium bag material. [0155] FIG. 17 to FIG. 20 show the design of the helium bag 20 and containment area for Version 2 of the design. This helium bag 20 is designed to have an ambient temperature volume of about 25% of the main envelope 1 . The fully expanded helium bag 20 would have a volume of about 85% of the main envelope 1 volume. For positive pitch axis pendulum stability it will likely be necessary to contain the helium to a central location as much as possible while it expands. This is accomplished using a high stretch silicone tubing or equivalent netting 22 . This netting 22 is attached to the pressurized air envelope 1 and it will tend to keep the helium in the central un-netted portion of the helium bag 20 as the helium expands yet the netting will stretch easily as it is required to when the helium increase slightly in pressure. [0156] FIG. 17 also introduces location of the furnace 23 in the air vehicle. The depicted furnace heats both the air chamber inside envelope 1 and the helium in the helium bag 20 . [0157] FIG. 17 and FIG. 19 show the helium bag 20 at about 25% of the volume of the main envelope. [0158] FIG. 18 and FIG. 20 show the helium bag 20 at about 85% of the volume of the main envelope. [0159] FIG. 21 to FIG. 27 depict a detailed design for a FIG. 17 furnace 23 that heats both the air chamber and the helium chamber. It is centrally located on the bottom of the main envelope 1 . Combustion air enters the furnace from the bottom, the heated air from the flames created in the FIG. 21 burner section 24 , is sucked up through a zig-zag path to the top of the FIG. 17 furnace 23 where it is forced to turn around and go back down through the central furnace flue. The combusted air then exits out the bottom of the furnace. While the hot combustion gasses zig-zag up the furnace, the zig-zag structure passes through the air and then the helium chambers inside the main envelope 1 . Circulation fans in both the air chamber and the helium chamber force air and helium though zig-zag chambers of the furnace to extract heat from then combusted gasses before they return down through the central furnace flue. [0160] FIG. 21 shows a cut-away view of the furnace with entry and exit gas flows depicted by arrows. This drawing also isolates the FIG. 21 burner section 24 , the FIG. 21 air section 25 , the FIG. 21 helium section 26 and the FIG. 21 helium fan section 27 . A main goal in designing this furnace was to keep the fan motors outside the main envelope 1 . [0161] FIG. 22 shows the FIG. 21 burner section 24 of the FIG. 17 furnace 23 . Beginning on the bottom is the main furnace fan motor 114 . The inlet disc 28 feeds air to the burner ring 116 . The combustion air fan 115 is attached to fan motor 114 . The main exhaust flue 29 diffuses the combustion air outside the fan motor 114 . All central exhaust flue rings are insulated. The combustion chamber cone 30 directs the hot air from the flames to the beginning of the zig-zag heat exchanger. The inlet plenum 32 directs the air from the air fans 31 into the beginning of the zig-zag heat exchanger. The top of the burner section begins with a section of zig-zag heat exchanger using a disc 33 , and rings 34 , 35 and 36 . [0162] FIG. 23 shows a section of the air heat exchanger. This section is repeated until the air heat exchanger is as tall as needed. It begins with rings 34 , 35 and 36 . Then a disc 33 is attached, followed by rings 34 , 35 and 36 again. Then it is topped off with a disc 37 . [0163] FIG. 24 shows a section of the helium heat exchanger. This section is repeated until the helium heat exchanger is as tall as needed. It begins with rings 34 , 35 and 36 . Then a disc 33 is attached, followed by rings 34 , 35 , 36 and 38 . Then it is topped off with a disc 39 . [0164] FIG. 25 shows the helium fan section of the FIG. 17 furnace 23 . It begins with rings 34 , 35 and 36 , followed by a disc 33 . Then by rings 34 , 38 and 40 . It is capped off with cap 41 and the helium fan 42 . [0165] FIG. 26 shows an x-ray view of the FIG. 17 furnace 23 . This allows the separator section 43 to be visualized. This section separates the air and helium heat exchangers of the FIG. 17 furnace 23 . [0166] FIG. 27 shows and exploded view of the FIG. 26 separator section 43 . Air exhaust ring 44 is followed by disc 45 and then helium entrance ring 44 . [0167] FIG. 28 shows the envelope cooling entrance ducts 46 in an open position on an isolated section of the rear of main envelope. When the air vehicle lands with its cargo there is a need to be able to reduce the buoyancy of the vehicle by the weight of the cargo very quickly. Since the weight of the cargo is lifted by the heated gasses in the envelope 1 , it is necessary to reduce the envelope temperature quickly so the cargo can be removed with the air vehicle floating away. If the air vehicle gondola/cargo carrier can be moored to the ground, then the cargo can be off loaded right away but it would still be preferable not to have to wait an hour or more to cool before being able to safely take off again. These cooling entrance ducts 46 , when opened, allow the propulsion unit 4 to go into reverse thrust and blow cool air through the envelope 1 . This air would come out similar flaps in the front of the airship. This would replace hot air with ambient air very quickly. The inlet flaps 48 are spring loaded shut but swing inward when the air pressure on the outside is greater than the pressure inside the envelope. The propulsion unit 4 in reverse thrust will have plenty of power to overcome the internal pressure and open the inlet flaps 48 . The cooling entrance ducts 46 are made of envelope fabric and the expanding triangles 47 are rigid material spring loaded to try to return to their closed positions as shown in FIG. 29 . When a cable running through the rear edge tubing 50 of the entrance ducts is pulled tight, the entrance duct will close as depicted on FIG. 29 . If the cable is released the cooling entrance ducts 46 will open up with any reverse wind flow. [0168] The inlet flaps 48 also have a round, over pressure relief valve 49 in them. This pressure relief valve 49 will let air out of the main envelope 1 any time the internal air pressure gets to high. [0169] FIG. 30 shows the forward over pressure valves 117 open which is what will happen when the FIG. 28 propulsion unit 4 goes into reverse with the FIG. 28 cooling entrance ducts 46 open. The over pressure valves 117 are also there to protect the main envelope 1 from over pressure. [0170] FIG. 31 shows Version 2 of the air vehicle with 4 engines and a cargo carrier 52 structure slung below the gondola. The gondola 51 is no longer a Cessna 172 fuselage. This Version 2 has an envelope 1 displacement of 24,000,000 cubic feet and is 1,075 feet long. These 4 propulsion engines are at least 2,000 horsepower each. There are two propulsion engines on the gondola 51 and one engine 53 on the nose of envelope 1 and one on the tail of envelope 1 . The engines on the nose and tail of envelope 1 may be fully or partially gimbaled to allow the engines to point at a significant angle left/right and/or up/down. With this feature the air vehicle could have full thruster positional control in hover mode. That is, the thrusters could move the vehicle left/right, forward/backward, up/down and yaw left/right. This would be very useful when hovering over a ship to re-fuel at sea. [0171] Also new in FIG. 31 is the permanently attached cargo carrier 52 slung below the gondola 51 . This cargo carrier allows loading and unloading cargo into a cargo hold that is bigger than the C5 Galaxy's cargo hold. This cargo carrier has a streamlined shape and is an air pressurized structure. [0172] FIG. 32 shows this cargo carrier 52 on the ground with its loading ramps 54 lowered and it front and rear access doors 55 open. When the doors are open and there is no air pressure in the cargo carrier 52 , the basic shape is maintained by a series of carbon fiber hoops. The cargo carrier is connected to the gondola 51 by two external cables 57 which are attached to four cables 56 . [0173] FIG. 33 shows the bottom of the cargo carrier 52 in the air. The loading ramps 54 seal against the lowered access door 55 (not visible). The flexible landing cushion membrane 58 is a much stronger and durable fabric then the upper fabric of the cargo carrier 52 . There will be a couple of feet of air that has to be compressed out of the area below the cargo carrier floor before the cargo carrier 52 settles on its FIG. 34 ball tires 60 under the cushion membrane 58 . [0174] FIG. 34 shows the cargo carrier 52 in the air with the cushion membrane 58 removed. This exposes the ball tires 60 that are the landing gear of the cargo carrier. Also shown is the pallet array 59 that makes up the cargo carrier 52 load carrying deck. This array is made up of thirty-three ten foot square pallets, each with a ball tire 60 in their centers. [0175] The cargo carrier 52 also contains a pressurized pilot cockpit. The floor of the cockpit 61 is visible here in FIG. 34 . [0176] FIG. 35 shows a bottom view of a single ten foot by ten foot by four foot pallet 62 in the pallet array 59 that makes up the cargo carriers load carrying deck. The basic structure consists of four welded aluminum trusses 64 . These trusses 64 are welded together the square shape shown. Into this square frame is welded the aluminum fuel tank (approximately 1,000 gallons) that has the ball tire 60 mounted on the bottom. The top surface where the loads are placed is made up of a composite honeycomb panel 63 . Each truss 64 has locating balls 66 mounted on one end and matching locating sockets 67 on the other end. This allows adjoining pallets to be keyed to each other. The socket array is then held together by tensioned cables running through the truss tubes. [0177] FIG. 36 is a section view of FIG. 35 pallet 62 showing the internal chambers. Chamber 71 is the fuel tank. Chamber 72 is the ball tire's 60 pressurized air container. The fuel tank cone 65 serves as a load transfer path to the ball tire 60 . The center post 70 also transfers deck load to the ball tire 60 . The fuel tank bottom 68 also transfer deck load to the ball tire 60 . The aluminum top of the fuel tank 69 is welded to the fuel tank cone and the FIG. 35 trusses 64 . [0178] FIG. 37 shows a few more items more clearly. To make the honeycomb panel 63 , the internal honeycomb 75 is faced with composites surfaces 74 , 76 , 77 and 78 . Pressure ring 73 is used to distribute ball tire mounting bolt loads and ensure an air tight seal. [0179] FIG. 38 to FIG. 40 show the details of the tensioned cables 79 that hold the pallet array together. Not only are the tensioned cables 79 under tension but they are also spring loaded. This allows some flexibility in the FIG. 34 pallet array 59 . This cargo carrier has the ability to land in rough terrain because of these sprung pre-tensioned cables 79 . [0180] FIG. 38 shows a section view of three trusses with the cables 79 holding them together. Truss 64 a has been raised up like the FIG. 34 ball tire 60 of the truss 64 a pallet is on a big rock. Rather than over stressing and bending the truss, the lower sprung cable 79 is stressed enough to compress the springs that are attached to it. This allows the truss 64 a to bend at the edge without breaking and the ball tire from truss 64 b touches the ground to relieve the load on the cables. [0181] FIG. 39 shows the details of the sprung end of the cable 79 . In one end of each large truss tube, a sleeve 82 is firmly welded in by a full end fillet weld 85 and a number of rosette welds 84 . This sleeve 82 transfers the full cable load to the large truss tube 83 . FIG. 35 Balls 66 and FIG. 35 sockets 67 are removable items and are snapped into place. A stack of spring washers 88 and two pressure washers 87 are slid into the large tube 83 to contact sleeve 82 . The end of the cable has a swaged on metal threaded rod 81 which is placed through the washer stack and a nut 86 is screwed onto threaded rod 81 . This nut is tightened until the cable has the required pre-tension. [0182] FIG. 40 shows a section view of the swaged cap end of cable 79 . The FIG. 35 ball 66 has been removed from this truss 64 a . Sleeve 82 has been welded at 84 and 85 as in FIG. 39 . This end of cable 79 has a cap end 80 swaged onto it. The cap end 80 mates into sleeve 82 and holds this end of the cable 79 securely in place. [0183] FIG. 41 shows Version 2 of the air vehicle on a single point ground mooring, in high winds that have shifted 45 degrees since the vehicle landed. Rather than having the gondola cables twist up when the wind shifts like this the gondola 51 can be made with a swivel in it and there is no need to reposition the cargo carrier 52 until the winds have shifted more than 200 degrees. This swivel is detailed in FIG. 43 . [0184] FIG. 43 shows the internal load transfer structure of the gondola and includes a swivel if needed. The cargo carrier cables 57 enter the bottom of the gondola through a slot in the floor. This slot is in a disc 99 that is allowed to rotate in a groove in the floor. The disc 99 has gear teeth on it to enable a gear motor 101 to rotate or lock the disc where desired. If the gear motor is moved out of gear meshing range the disc would rotate freely, due to the swivel pin 103 , if the wind changed direction and the gondola cables 7 began to rotate with the wind change. If you were to takeoff with the cargo carrier out of line with the main FIG. 1 envelope 1 you could engage the gear motor 101 and straighten the FIG. 41 cargo carrier 52 out. Pulley 97 has a narrow angle groove in it so that when the gondola cable 7 is loaded, the gondola cable 7 is tightly wedged into the pulley and it will not slip, however gear motor 102 could still rotate the pulley and raise or lower the nose of the FIG. 1 envelope 1 . Frame 100 attaches to the roof of the FIG. 41 gondola 51 and carries the whole weight and thrust loads of the FIG. 41 gondola 51 . [0185] The FIG. 8 gondola cables 6 and 7 and the FIG. 32 cargo carrier cables 57 and 56 present a significant drag penalty due to their frontal area, shape and length. If the cables were streamlined the amount of drag becomes insignificant. Therefore FIG. 42 shows a collapsible streamlined cuff that could be used on section of the cables that collapse and expand ( FIG. 43 cable 7 near the gondola). Also, a lot of the cables will have fuel or electric lines running up and down them. The cuff presents a way to deal with those fuel and other lines by running them in chambers 94 or 95 of the cuff. [0186] In FIG. 42 , cuff 89 R is about to be joined to cuff 89 L with screws 96 . When this new 89 LR cuff is screwed together so that lip 92 of cuff 90 c is inside and above lip 91 of cuff 89 L, the new 89 LR cuff should slid up and down on cuff 90 c . You can also see the collapsing ability by looking at cuff 90 a and 90 b which are fully collapsed. Cuff 90 b and cuff 90 c are fully extended. The heavy load cable runs in the round chamber 93 . On sections of the cables that do not need to extend and contract, a cuff similar to cuffs 89 L and 89 R can be made with without the extending feature. [0187] FIG. 44 shows a front skeleton view of the FIG. 51 cargo carrier 52 with the fabric cover removed, the loading ramps 54 down and the front access door 106 open. There are seventeen carbon fiber hoops 104 that hold the removed fabrics shaped when the access doors 106 open and the shell looses pressure. You can clearly see the pallet array 59 . The cockpit 105 is raised and lowered with the access door 106 but it is still accessible by ladder. In the rear of the FIG. 51 cargo carrier 52 and on the tray 107 , are two electric generators 108 used for providing electric power to all areas of the vehicle. Only one generator is needed at any time and the other is a backup. This is also the generator that needs to run when to keep the FIG. 1 envelope 1 pressurized while moored on the ground. FIG. 45 is a rear view of the same skeleton with the front access door closed and the hoops 104 removed. Exposed now are the internal cargo carrier load cables 118 , these load cables further distribute the cargo carrier cable loads to the pallet array 59 . [0188] FIG. 46 to FIG. 51 depict the procedure to do a full maintenance mooring at a site 109 equipped with the nine mooring points that allow full maintenance of the vehicle [0189] FIG. 46 shows Version 2 of the air vehicle having just landed at the site 109 near the single point mooring location 111 . At this point a mooring cable is extended from a winch on the FIG. 51 cargo carrier 52 and attached to the single point mooring location 111 . The pilot then takes off again as in FIG. 47 and winches the air vehicle back to the ground, but this time the vehicle is position precisely over the cargo carrier mooring points. [0190] Once the FIG. 51 cargo carrier 52 is moored down securely, the four maintenance FIG. 48 mooring lines 16 are commanded to extend from small winches at the FIG. 3 load bridle 5 and FIG. 3 gondola cable 6 junctions. These mooring cables 16 are then attached to the four stronger ground mooring winches as depicted in FIG. 48 . [0191] The ground winches begin winching the main FIG. 1 envelope 1 to the ground. As it lowers, the FIG. 41 gondola 51 will also be lowered to the ground where it will settle into a maintenance FIG. 49 saddle 113 . [0192] FIG. 50 shows the 1075 foot long vehicle in its full maintenance position. At this point, short lines can replace the ground winches so the winches can be used elsewhere. [0193] FIG. 51 shows a joystick controller 119 that is used to control any number of thrusters attached to the air vehicle. [0194] When the joystick 120 is maneuvered in any of the indicated directions, the appropriate thrusters will act to add thruster force to the vehicle that would result in movement in the indicated direction if the vehicle were not being constrained. The amplitude of the added thruster force will be proportional to the maneuvering effort on the joystick 120 . In the case of a variable pitch propeller thruster, first the pitch would be increased to its maximum and then the engine power would be increased to its maximum. [0195] The word maneuvering is used to indicate either a displacement in the indicated direction or a force is applied in the indicated direction on the joystick 120 . Any of the sensors in the joystick control 119 that sense this maneuvering effort on the joystick 120 could be either a force sensor, a displacement sensor or both a force and a displacement sensor. [0196] Any of the maneuvering directions may be allowed to return to a neutral condition when the maneuvering effort is removed from it. In this case, the direction will also have a trim slide that can change the commanded signal at the no maneuvering effort output command position. [0197] While there have been shown and described and pointed out the fundamental novel features of the invention as applied to the preferred embodiments, it will be understood that the foregoing is considered as illustrative only of the principles of the invention and not intended to be exhaustive or to limit the invention to the precise forms disclosed. Obvious modifications or variations are possible in light of the above teachings. The embodiments discussed were chosen and described to provide the best illustration of the principles of the invention and its practical application to enable one of ordinary skill in the art to utilize the invention in various embodiments and with various modifications as are suited to the particular use contemplated All such modifications and variations are within the scope of the invention as determined by the appended claims when interpreted in accordance with the breadth to which they are entitled.

A system for controlling the lift of an airship for carrying a cargo and a supply of fuel has a self supporting hull made of a flexible gas impermeable material. Air which is located in the hull is at an internal pressure which is greater than the atmospheric pressure of the air on the outside of the hull, and a bag filled with helium is located within the hull and is surrounded by the air in the hull. A heating means is provided to heat both the air and the helium. The bag located within the hull and surrounded by the air in the hull has enough helium in it at ambient temperature to lift all but the cargo and fuel in the airship, and the air and helium, when heated, provides increased buoyancy due to the increase of expansion of both heated gases to lift the airship with its fuel and cargo.

BACKGROUND OF THE INVENTION 1. Field of the Invention This invention relates to an optical transmission element comprising a fibrous light waveguide surrounded by a cushion layer with both positioned within a sheath of a hard material. There is further provided a glide layer between the sheath and the cushion layer for accommodating relative movement between the sheath and the cushion layer and waveguide. 2. Description of the Prior Art Optical transmission elements which contain a fibrous light waveguide surrounded by a cushion layer are known in the art. See, for example, DE-OS No. 30 11 009. The element disclosed therein is surrounded by a relatively stiff protective sheath. A parting layer is provided between the cushion layer and the stiff protective layer, and the parting layer permits movement of the light waveguide and the cushion layer on one hand and the outer protective sheath on the other hand. In use and during processing, such as stranding and the like, the light waveguide may be subjected to undesirable mechanical stresses, and it is desirable to protect the waveguide from such stresses. It is therefore an object of this invention to provide an improved hard protective sheath which maximizes protection of the sensitive optical fibers of the light waveguide, particularly during various processing procedures. Another object of this invention is to minimize surface and stress cracking which may occur in the hard protective sheath. SUMMARY OF THE INVENTION There is provided by this invention an optical transmission element of the type described above in which protection of the optical fibers is maximized by said sheath having an elastic modulus (E) greater than 2000 N/mm 2 and a coefficient of thermal expansion less than 0.8×10 -4 /K, preferably between 0.5×10 -4 /K and 0.7×10 -4 /K. The element also includes a glide layer between the cushion layer and hard sheath for enhancing protection of the light waveguide. When needed, an outer protective layer is provided over the hard sheath for minimizing surface and stress cracking. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is a cross-sectional view of an optical transmission element having a light waveguide, cushion layer, glide layer, hard outer sheath and outer protective layer. DESCRIPTION OF THE PREFERRED EMBODIMENT Referring first to the drawing, there is shown an optical transmission element (AD) which includes a central light waveguide (LW), surrounded by a cushion protective layer (PS), a glide layer (GS), a hard outer sheath (HL), and a outer protective layer (SU). Referring now to the hard outer sheath (HL), the use of materials having an elastic modulus (E) greater than 2000 N/mm 2 assures a sufficiently stiff sheath. Furthermore, due to the sheath having such an elastic modulus, (1) deformations of the sheath due to local shearing forces (E·A≃1000 N) remain slight and (2) shearing forces (E·A) in the longitudinal and cross-sectional direction of the glass fibers (which are approximately 125 μm in diameter) and the sheath are of approximately the same magnitude. This results, on one hand, in approximately equal division of all forces (roughly 1:1) and, on the other hand, minimizes the risk that the higher coefficient of expansion of the sheath will lead to deformations injurious to the fibers. For the sheath, coefficients of expansion below 0.8×10 -4 /K are preferred because the difference in coefficients of expansion of the fiber and the sheath remains within about one order of magnitude. Assuming a ±50° K. temperature deviation, the maximum forces of ##EQU1## act on the sheath and the fiber according to the equation ##EQU2## These forces are to be viewed as acceptable both mechanically as well as optically. The thin glide layer (GS) also is temperature responsive, and thus provides a supporting function for the light waveguide fiber (LW) at low temperatures and resists buckling by the crushing forces of this magnitude. The transmission element constructed in accordance with the invention offers protection for the light waveguide (LW), which approximates that of a light waveguide loosely carried within a tubular sheath. An element constructed in accordance with the invention can thus advantageously be stranded into a bundle together with other similarly constructed transmission elements, even in a tightly packed core format. This could not be achieved without appropriate dimensioning of the sheath nor without the high degree of mechanical decoupling between the outer sheath and the cushioned light wave resulting from the use of the glide layer (GS). The glide layer (GS) is of significance in this context and is fabricated from a cross-linked or thixotropic filling compounds which have a bearing or load carrying ability. The elastic modulus values of the glide layer are below 0.01 N/mm 2 and values between 0.01 N/mm 2 and 0.001 N/mm 2 are particularly suitable. The elastic modulus for the individual components are roughly as follows: E modulus LWG glass (LW): 78000 N/mm 2 E modulus coating (PS): 10 N/mm 2 E modulus filling compound (GS): 0.01 N/mm 2 E modulus sheath (HS): 2000 N/mm 2 The thickness of the glide layer (GS) is selected between 50 and 100 μm, so as to cooperate in permitting relative movement between the hard sheath (HL) on one hand and the cushioned light waveguide fiber (LW) on the other hand. Movement of the sheath (HL) is taken up or compensated for by the glide layer (GS) and therefore does not reach the light waveguide (LW). The cross-sectional areas for the various elements are approximately as follows: Light waveguide fiber (LW): 1/80 mm 2 Cushion layer (PS): 0.2 mm 2 Filling compound (GS): 0.2 mm 2 Sheath (HL): 0.3-0.4 mm 2 Based on these cross-sectional areas, the local shearing forces (E×A) are as follows: Light waveguide fiber (LW): 1000 N Coating (PS): 2 N Filling compound (GS): 0.2 N Sheath (HL): 600-1500 N The cushion layer (PS) is desirably of an ultraviolet cross-linked material since at high production speeds freshly drawn fiber can be coated at high speeds using such materials. Moreover, if necessary, removal of the cushion layer (PS) is possible using simple swelling agents, or under certain conditions solely mechanical removal is possible. By comparison, silicone rubber adheres very firmly to the fiber. Since the outer sheath (HL) should be relatively stiff, materials that are hard and tough are principally considered, that is, raw materials for injection molding of the type which produce high impact resistant housings whose tenacity prevents crack formation. Such injection molding materials includes aramides and polyesters. Other materials which could be suitable include certain polycarbonates, but due to their high hardness, they are relatively susceptible to stress crackings. This means that fine surface cracks can appear due to the deformations resulting from bending stresses, and water can penetrate through the cracks into the element. The reason for such stress cracking relates to the high frozen-in orientation stresses that can spread and deepen due to surface attacks by ultraviolet light or solvents, even for example, from the glide layer (GS). Since the substances used herein contain a broad mixture of molecular weights --due to manufacturing considerations or intentionally selected to obtain a specific melt index--the stress crack behavior increases with increasing hardness and elastic modulus. The stress cracking of the sheath can be minimized by coating or applying a protective layer to the exterior surface of the hard sheath, which layer is not or is only very slightly susceptible to stress cracking as compared with the sheath. High molecular weight materials on a polyester base are preferably employed for the protective layer itself. A preferred combination for the sheath is grilamide/polybutylenterephthalat which is not susceptible to stress cracking either as a combination or individually. However, the aramide can absorb up to 3% moisture and will thus change in dimension depending upon moisture content. The use of a moisture-inhibiting and moisture-resistant polyester layer as an additional protective layer minimizes this potential swelling and cracking problem. The ultraviolet sensitivity of the aramide layer is also eliminated by means of the polyester protective coating. On the other hand, the polyester layer alone would not be hard or resistant enough to swelling agents (such as found in the glide layer), even if it had a greater thickness, so that the preferred combination of the coated hard sheath offers an optimum stability and protective effect. The transmission element constructed in accordance with the invention is employable in a particularly advantageous fashion as an individual lead for the formation of stranded lead bundles since the mechanical deformation forces occurring in stranding cannot penetrate or at least cannot penetrate disruptively to the light waveguide due to the hard protective sheath and due to the intervening glide layer. In view of the use of the transmission element as a lead for the core of an optical cable, one selects the material of the protective layer such that swelling by the potential core filling compound is minimized. Based on the foregoing general discussion, reference is again made to the drawing. From the center toward the outside, the element includes the light waveguide (LW) which is surrounded by the protective or cushion layer (PS). Ultraviolet cross-linkable materials are preferably employed for the cushion layer (PS). The cushion layer rests directly on the outer surface of the light waveguide (LW). The cushion layer (PS) is surrounded by a glide layer (GS) which permits movement between the hard outer sheath (HL) and the cushion layer (PS). It is desirable that the glide layer (GS) have minimal bearing properties, and preferably it consists of easily cross-linked or thixotropic material having a low intrinsic viscosity. In order, however, to further assure the decoupling or separation of the light waveguide from movements of the sheath, the glide layer (GS) may only have a very low elastic modulus, advantageously below 0.01 N/mm 2 , preferably between 0.01 and 0.001 N/mm 2 . The outside diameters of the individual component parts of the transmission element are approximately as follows: Light waveguide (LW): 50-125 μm Cushion layer (PS): 250-500 μm Glide layer (GS): 400-700 μm Sheath (HL): 700-1000 μm In order to protect the light waveguide fiber (LW) which is surrounded by a soft and thus mobile and cushioned manner from external mechanical stresses, the sheath (HL) is fabricated from a correspondingly hard material such as aramides, polyethersulfones, polycarbonates, etc., being specifically employable for this purpose. Some such hard materials have the disadvantage that orientation stresses exist to a high degree and lead to superficial cracking of the hard layer or sheath due to exposure to ultraviolet light, solvents or wetting agents. This situation as described above can be minimized with a protective layer (SU) which is not susceptible to stress cracking and which is applied to the outside of the hard layer (HL). The protective layer is so elastic that any and all cracks on the inside layer or (HL) are minimized. The outer protective layer (SU) desirably consists of a highly plastic material on a polyester base which can be applied by means of a tubular die immediately following the extrusion process for the stiff sheath (HL). One such material is sold under the name Vestodur B 3000 (polybutylene terephthalate) by Firma Chemische Werke Huels (4370 Marl, West-Germany). The wall/thickness of this protective layer (SU) can be relatively thin, between 10 and 80 μm and advantageous at 50 μm. It should amount to between 10 and 50%, preferably about 30% of the wall thickness of the sheath (HL). Even relatively thin wall thicknesses of the protective layers already form adequate protection against stress cracking of the stiff sheath (HL). A protective layer (SU) is applied to the sheath (HL), such that it adheres to the latter, so that no lifting or separation of layers from one to another occurs in the event of bending, which possibly could lead to buckling under severe bending. In the ideal sheath combination, moreover, the protective layer has strengths which complement the weakness of the individual layers, for example, water sensitivity of aramides, oil sensitivity of the polyesters, etc. The wall thickness of the glide layer (GS) advantageously lies between 50 μm and 100 μm; the wall thickness of the sheath (HL) between 80 μm and 140 μm, whereby the values of the diameter from 700 μm to 1100 μm are expedient. The aforementioned dimensions are to be applied in a particularly advantageous fashion given light waveguide leads that are stranded into bundles. Although the invention has been described with respect to preferred embodiments, it is not to be so limited as changes and modifications can be made which are within the full intended scope of the invention as defined by the appended claims.

The optical transmission element (AD) comprises a light waveguide (LW) provided with a cushion layer (PS), said light waveguide being positioned within a hard sheath (HL). A glide layer (GS) is provided between said sheath (HL) and said cushion layer (PS). The sheath (HL) is of a material having a modulus of elasticity greater than 2000 N/mm 2 and a coefficient of thermal expansion less than 0.8×10 -4 /°K.

CROSS-REFERENCE TO RELATED APPLICATION This application is a continuation of U.S. patent application Ser. No. 09/821,182, filed Mar. 29, 2001, now U.S. Pat. No. RE41,733, which is a reissue of U.S. patent application Ser. No. 08/863,156, filed May 27, 1997, now U.S. Pat. No. 5,889,694, which is a continuation-in-part of U.S. patent Application Ser. No. 08/610,992 entitled “Dual Addressed Rectifier Storage Device”, filed Mar. 5, 1996, now U.S. Pat. No. 5,673,218, issued Sep 30, 1997. Reissue application Ser. No. 11/780,300, filed Jul. 19, 2007 is also a resissue of U.S. Pat. No. 5,889,694. The entire disclosure of each of these applications is hereby incorporated by reference. BACKGROUND OF THE INVENTION 1. Field of the Invention The present invention relates to electronic data retrieval devices, and more particularly to electronic digital logic devices having semiconductor mass storage capabilities by virtue of their data being stored in highly symmetrical arrays of diodes. 2. Prior Art Most present day devices having mass storage capabilities rely on such moveable media as magnetic disks, optical compact disks, digital tape or the like. Some devices having large storage Capabilities have utilized large numbers of read-only memory (ROM) devices. Many read-only memory (ROM) devices have been disclosed having a wide variety of implementations. In many of these devices the bit storage means is accomplished through the application of gates or transistors. But, a subset of these ROM devices has accomplished the bit, storage means through the use of a matrix of diodes, as was disclosed by Robb in U.S. Pat. No. 3,245,051. Many of these ROM device is include diode matrix storage means utilizing a set of conductors that act as selectors and a second orthogonal set of conductors that act as data outputs. In U.S. Pat. No. 4,070,654, one set of generally parallel conductors acts as the Selection Input Lines and a second set of generally parallel conductors that is orthogonal to and overlapping with the first acts as the Digit Output Lines. A bit of information is represented at each point of intersection of the Selection Input lines with the Digit Output Lines by the presence or absence of a diode at that point, where the presence or absence of a diode distinguishes the logical state of the stored information bit at that point of intersection. A selection circuit selects one line of the Selection Input Lines such that the state of all of the Digit Output Lines is then controlled to the extent that each of those Digit Output Lines is connected to that selected Selection Input Line through a diode. All of the digit Output Lines are read in parallel. A disadvantage is that as the matrix is increased in size, the complexity of the selection logic that drives the selection circuits (such that one line is selected out of the many Selection Input Lines) grows exponentially. But as this matrix increases in size, so too will the number of Digit Output Lines that will have to be simultaneously supplied with current and that current will vary depending upon the state of the bits at those various locations. Also, as the number of simultaneously driven Digit Output Lines increases, some means of selecting the subset of desired data bits would have to be added. In U.S. Pat. No. 4,661,927, some of the problems of the exponential growth in the complexity of the addressing circuitry and of the loading on the selected addressing lines are dealt with. Addressing is accomplished by using diode-transistor logic (DTL) for the input addressing. The transistors of the DTL selection circuits act as buffer-drivers between the address selection means and the bit storage means thereby providing the current needed to source a growing number of data bit output lines. However, the transistors in this DTL circuitry add complexity to the overall circuit which will reduce packing densities and add an additional problem—that of leakage currents in those transistors—that requires additional compensating circuitry that further reduces packing densities. The use of dummy diode loads for balancing data dependent loading variations reduces packing densities even further. In U.S. Pat. No. 4,884,238, the problem of loading is dealt with by utilizing FET switches to disconnect all but the desired Digit Output line. The selected bit is present at the intersection of two selected orthogonal conducting lines. In this way, the number of bits simultaneously selected does not grow with the size of the array and the problem of loading can be controlled. But, such a design still requires a large number of FET transistors and the addressing means to control those FET transistors and the interconnection wiring to connect said addressing means with said FET transistors which will reduce packing densities. While the addressing means could be of the same DTL type to keep said addressing means small, the large number of buffering transistors to select the various conductive lines will remain large (at least one FET per conductive line and, on about half of the lines, two FET's). Furthermore, this inclusion of FET transistors may make the device subject to damage from static electrical discharges that might make it less practical for use in a consumer product where the consumer may handle these devices. In U.S. Pat. No. 4,347,585, Eardley discloses a dual-addressed device wherein the selected diode is at the intersection of a column line and a row line (where the column lines are connected to the diodes' cathodes and the row lines are connected to the diodes' anodes) such that the voltage potential of one column line is lowered and the voltage potential of one row line is raised thereby forward biasing the diode, if any, at the point of intersection of the two lines. One of the features of this device over the prior art (as discussed in that patent) is the circuitry for the selection logic; Eardley discloses means for line selection comprising high speed transistor driver circuits. The disclosed device also requires two types of Schottky barrier diode devices. As a result, it is anticipated that the disclosed device will suffer from several problems, particularly when one attempts to scale up the device to extremely high storage densities. These problems may include transistor current leakage becoming noticeable as the number of transistors increases and device yields becoming reduced as the complexity of multiple semiconductor fabrication steps and nore complex device interconnect circuitry increases. These problems may prevent the kind of size scaling that could result in devices in the Gigabit range that would be necessary to create memory chips that could replace today's CD-ROMs. As will be shown below, the Dual-addressed Rectifier Storage DRS) Array comprised by the present invention solves many of the Problems associated with the above mentioned inventions while staining high data packing densities by simultaneously using both orthogonal sets of conductors to address the data bits without the need for transistor switches on each conductive line. It does this by having a diode-logic addressing mechanism directly controlling the voltage levels on the conductive lines. Also, by extending the application of the diode array to perform the functions of addressing, storage, and bit sensing, symmetry is increased and this higher symmetry results in higher packing densities. In U.S. Pat. No. 4,070,654, among others, a means is disclosed for programming the information into a semiconductor diode array by selectively etching away openings through the oxide layer that insulates the plurality of doped conductors from the orthogonal plurality of metalized conductors on the surface such that each opening enabled contact between the respective conductor of each plurality thereby forming a diode representing a toggled bit of stored information at that array location. The present invention discloses a means for constructing the semiconductor device up to the final metalization etch step before programming the data thereby enabling the programming of data to be performed much later in the manufacturing process. Mass storage devices comprising moveable media such as magnetic disks, optical compact disks, digital tape, or the like, have motors and other mechanical parts that are prone to breaking or wearing out, can suffer audio disruption when subjected to vibrations, are too heavy to be carried during certain activities such as jogging, and consume significant electrical power (due to the operation of the mechanical components). Devices utilizing ROM chips are limited in their capacities due to the limited storage densities of present day ROM chips. The present invention eliminates or reduces all of these drawbacks because it uses DRS Arrays and, as a result, has the high storage densities of a CD-ROM without having mechanical parts. SUMMARY OF THE INVENTION Instead of a CD-ROM and its associated mechanical components, this device comprises one or more Dual-addressed Rectifier Storage (DRS) Arrays which are read-only memory (ROM) devices that utilize an array of rectifiers for its storage means. Like the predecessors to the DRS Array, the logical state of stored data is determined by the presence or absence of rectifiers at the points of intersection of two orthogonal and overlapping sets of generally parallel conductive lines. The present invention uses both sets of generally parallel conductive lines for addressing but senses the logical state of the addressed data by sensing the loading on the selected lines. The DRS Array comprises a cross-point selection means that enables the selection of a single point of intersection from within an array by applying a forward voltage across that point; this selection means will find applications in Read Only Memory (ROM) as shown in the present device, One-Time Programmable Read Only Memory (OTPROM), Random Access Memory (RAM), and LED matrix displays. More specifically, the DRS Array comprises an array of Rectifiers (where the column lines are connected to the cathodes of said rectifiers and the row lines are connected to the anodes of said rectifiers) wherein the row lines and column lines are pulled through resistive means to either the positive supply or to ground, respectively, such that, absent any addressing circuitry, all of the rectifiers in the array would be forward biased. The addition of addressing circuitry will selectively connect those row lines and column lines to either ground or the positive supply such that the voltage potential between any row line and column line would be dropped to the point that an interconnected rectifier no longer would be forward biased. Any line whose voltage is pulled close to either ground or to the positive supply and away from that voltage that would result from the resistive means alone will be referred to as Being “disabled.” Selection of a line in both sets of generally parallel lines is accomplished by a diode addressing array similar to that disclosed in U.S. Pat. No. 4,661,927 but, instead of using diode-transistor logic (DTL) means to do this addressing, no transistor buffer-driver stage is used, thereby greatly reducing circuit complexity while eliminating a source of possible current leaks. This is possible by taking advantage of the forward voltage drop characteristics of a rectifier to control the voltage levels. The present device does not attempt to drive a heavily loaded selected line (where all of the orthogonal output lines that are connected to the selected line through diodes comprise that loading) that might require such a buffer-driver stage; the only load is that rectifier connection to the selected orthogonal line. One result of this disabling means is a greatly simplified selection circuit on both orthogonal sets in the array. The addressing means, the storage means, and the bit sensing means are all formed in the same rectifier array structure in a highly simple and symmetric design that is ideal for high packing densities. The present invention comprises DRS Arrays that could be removable and interchangeable so that one could pop them in and out according to their current musical interest or desire. A DRS Array is expected to be able to hold the equivalent of an entire CD-ROM or more on a single one inch (or smaller) square of silicon, unlike conventional read only memory chips (ROM's) which would require about 50 ROM chip if each contained about 10 Megabytes of data. By contrast, the present invention would be very small and compact. While most present day mass storage devices rely on such moveable media as magnetic disks, optical compact disks, digital tape or the like, the present invention, through the use of the DRS array, eliminates the mechanical components of such devices. In so doing, the present invention reduces the risk of mechanical failures, the problems of bulkiness, and the consumption of electrical power. The many possible variations on the DRS Array make it a very versatile storage medium ranging from a stand-alone array chip, to an array chip with an incorporated sequentially loaded address sequencer, to a microcomputer chip that utilizes the DRS Array for its program memory, to a personal stereo unit utilizing a removable DRS Array instead of a music CD-ROM, to a pocket sized video player utilizing a removable DRS Array module containing compressed video data. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 . illustrates a block diagram of a digital logic device comprising a Dual-addressed Rectifier Storage (DRS) Array, keypad, CD display and dual analog outputs. FIG. 2 . illustrates a schematic diagram showing one way to interface a Dual-addressed Rectifier Storage Array to a digital logic device. FIG. 3 . illustrates a schematic diagram of a Dual-addressed Rectifier Storage Array. FIG. 4 . illustrates a schematic diagram of a Dual-addressed Rectifier Storage Array which includes complementary address selection circuitry. FIG. 5 . illustrates a variation on the data sensing means of a Dual-addressed Rectifier Storage Array for simultaneously accessing multiple stored bits in parallel and a variation on the addressing means for reducing the device's power consumption. FIG. 6 . illustrates the doping step in the semiconductor manufacture of the Anode Lines, Cathode Lines, and rectifiers comprised by a Dual-addressed Rectifier Storage Array. FIG. 7 . illustrates the oxide growth step in the semiconductor manufacture of the Anode Lines, Cathode Lines, and rectifiers comprised by a Dual-addressed Rectifier Storage Array. FIG. 8 . illustrates the oxide etch step in the semiconductor Manufacture of the Anode Lines, Cathode Lines, and rectifiers comprised by a Dual-addressed Rectifier Storage Array. FIG. 9 . illustrates the metalization step in the semiconductor Manufacture of the Anode Lines, Cathode Lines, and rectifiers comprised by a Dual-addressed Rectifier Storage Array. FIG. 10 . illustrates the metalization etching step in the semiconductor manufacture of the Anode Lines, Cathode Lines, and rectifiers comprised by a Dual-addressed Rectifier Storage Array. FIG. 11 . illustrates a variation on the semiconductor manufacture of the Anode Lines, Cathode Lines, and rectifiers comprised by a Dual-addressed Rectifier Storage Array. FIG. 12 . illustrates a plot of the voltage/current relationship of a diode. FIG. 13 . illustrates a schematic diagram of a variation on the Dual-addressed Rectifier Storage Array wherein the resistive means is accomplished via “leaky” diodes. FIG. 14 . illustrates a schematic diagram of a variation on the Dual-addressed Rectifier Storage Array which includes serial addressing circuitry. DESCRIPTION OF THE PREFERRED EMBODIMENT Refer now to the drawings which show a preferred embodiment of the invention. FIG. 1 shows a block diagram of a digital logic device comprising a microcomputer, a keypad, an LCD display, analog to digital converters, analog buffers and amplifiers, a headphone jack, and a Dual-addressed Rectifier Storage (DRS) Array. Microcomputer, A, is connected to a keypad, B, and an LCD display, C. This configuration is very common and several of the manufacturers of microcomputer chips have application notes showing the details and schematics of such a circuit configuration. Microcomputer, A, is also interfaced to two 16 bit digital to analog converters, E, whose outputs are connected to circuitry, F, capable of driving a set of headphones (plugged into this device at the headphone jack, G). Circuitry such as this just described is known by one skilled in the art; this circuitry exists in essentially this form in many common devices, including portable CD-ROM audio players used for listening to music CD's. Software running in the microcomputer, A, would cause that microcomputer to sequentially read 4 bytes of data from the DRS Array, D, every 25 μSec and move that data to the two 16-bit digital to analog converters, E. The result of this process would be 16 bit stereo sampled at 40 KHz thereby matching the audio of current CD-ROM technology. Software running in the microcomputer could also cause information read from the DRS Array (such as musical artist, title of the audio tracks, playing time, and the like) to be displayed on the LCD display. Software running in the microcomputer could also read from the keypad to affect the operation of that running software and cause it to jump to a specific address in the DRS Array at which point it would continue to read sequentially (such as to skip or repeat certain audio tracks, play the tracks contained in the DRS Array in random order, or the like). The data contained in the DRS Array need not be limited to musical data. Video data, computer software or applications, reference data such as text, diagrams or the like, or a variety of other information could be stored. Of course, in many of these possible data types, the exact configuration as shown in FIG. 1 may not be needed. for example, a DRS Array containing computer software would likely not need two digital to analog converters for output but would rather need interface logic such that it could be connected to a standard computer such that it emulated a standard CD-ROM drive. In this way, one could enjoy the benefits afforded by the use of DRS arrays without having to modify one's standard computer or its software (given an accurate emulation of a CD-ROM drive, that standard computer would operate just as if it was actually connected to a CD-ROM drive). What is needed is the ability to connect DRS Arrays to such digital logic devices as a personal computer, a microprocessor, a microcomputer chip or the like. Refer now to FIG. 2 which shows a possible way to interface a Dual-addressed Rectifier Storage Array to a digital logic device. On microprocessor, A, the address lines (or input/output ports configured to perform the function of addressing), A 0 through A 15 , connect directly to the lower address lines, A 0 through A 15 , of the DRS Array, as well as to the inputs of the 16 bit latch, H. The upper 16 bits of the address to the DRS Array are set by writing to a location in memory having the same lower 16 address bits as is desired for the upper address bits. When the write line, −WR, goes low and, at the end of that write cycle (when that write line goes back high), that rising edge will cause the address bits A 0 through A 15 to be latched in the 16 bit latch, H, where they will be held to address the upper address lines of the DRS Array, A 16 through A 31 , until the next write cycle. Note that the data byte written is ignored, only the address bits are needed to set the latch. (A variation would be to make use of that data byte to control the selection of multiple DRS Arrays or to address additional address bits of an even larger DRS Array. Other techniques for latching additional selection bits are readily known by those skilled in the art if multiple DRS Arrays are to be interfaced. Of course, if the digital logic device was a microprocessor having an addressing port sufficiently wider than the 16 bits shown in this figure, then the latching mechanism would be unnecessary as the entire DRS Array might be addressed directly.) The DRS Array is configured to be enabled by applying power through PNP transistor, I, and resistor, J, when the read line, −RD, goes low. The DRS Array as shown in FIG. 2 would contain about 537 Megabytes of information or about 56 minutes of 16-bit stereo audio sampled at 40,000 KHz, roughly the equivalent of a present day CD-ROM. Throughout the remainder of this description, references will made to the positions of different parts of the DRS Array; this is for ease of looking at the figures only, and is not meant to imply physical location of components in an actual device. Also in this description, transistors in the fully on state will be referred to as being saturated. Switching transistors into and out of saturation is typically slower than switching them into and out of a nearly saturated state (although nearly saturated will generally work as well). Circuitry to keep transistors from becoming fully saturated is desirable but not required and has been omitted for the purpose of keeping the descriptions and figures uncluttered. Refer now to FIG. 3 which illustrates an example of a preferred embodiment of a Dual-addressed Rectifier Storage Array. While the array in this example is only 64 bits in size, it will become clear to one skilled in the art that this array is highly scaleable. Rows are drawn running horizontally in this figure and Columns are drawn running vertically. The Storage Rows, P, are connected to the Row Resistors, R. Also connecting to the Storage Rows, P, at various points are the anodes of the Storage Rectifiers, K, the anodes of the Row Addressing Rectifiers, L, and the anodes of the Storage Bit Sensing Rectifiers, U. Connecting to the Addressing Columns, N, are the cathodes of the Row Addressing Rectifiers, L. Connecting to the Storage Bit Sensing Column, M, are the cathodes of the Storage Bit Sensing Rectifiers, U. The Storage Columns, O, are connected to the Column Resistors, S. Also connecting to the Storage Columns at various points are the cathodes of the Storage Rectifiers, K, and the cathodes of the Column Addressing Rectifiers, T. Connecting to the Addressing Rows, Q, at various points are the anodes of the Column Addressing Rectifiers, T. To understand the operation of the device, first consider what the operation of the device would be absent the row addressing means (the Row Addressing Rectifiers, L, and the Addressing columns, N) and the column addressing means (the Column addressing Rectifier, T, and the Addressing Rows, Q) and the bit sensing means (the Storage Bit Sensing Rectifiers, U, and the storage Bit Sensing Column, M). What would be left are the Storage Rows, P, which are each pulled to the positive supply through the Row Resistors, R, and the Storage Columns, O, which are each pulled to ground through the Column Resistors, S. The result would be that any rectifier present among the Storage Rectifiers, K, would be forward biased and the forward voltage drop of any said Storage rectifier would be centered around the voltage level of one-half of the positive supply (the Row Resistors, R, and the Column Resistors, S, are of equal resistance values and therefore form a voltage divider having a center voltage level of one-half of the positive supply). The resistive means, Row Resistors, R, and the Column Resistors, S, could be constructed through the use of resistors or their equivalent, the use of transistors (bipolar or FET) biased in their linear region (between being completely turned off and being saturated), the use of “leaky” diodes, or the like. Next, consider the impact of the Row Addressing Rectifiers, L, and the Addressing Columns, N. The Addressing Columns, N, work in complementary pairs labeled by address designation An and its complement −An (where n indicates the Nth address line). Addressing is performed by pulling an Addressing Column near to ground or near to the positive supply and its complementary Addressing Column near to the positive supply or near to ground respectively. (Allowing an Addressing Column to float instead of pulling it close to the positive supply would work, too.) When an Addressing Column is pulled near to ground, any of the Row Addressing Rectifiers whose cathodes are connected to said Addressing Column will be forward biased and the voltage drop across those rectifiers will determine the resulting voltage on the Storage Rows connected to said forward biased Row Addressing Rectifiers. Storage Rows whose voltages are pulled down in this way are called “disabled.” The voltage level on said disabled Storage Rows would be equal to the near to ground voltage on the Addressing Column plus the forward voltage across said forward biased Row Addressing Rectifier. Next, consider the impact of adding the Column Addressing rectifiers, T, and the Addressing Rows, Q. The Addressing Rows, Q, work in complementary pairs labeled by address designation An and its complement −An (where n indicates the Nth address line). addressing is performed by pulling an Addressing Row near to ground or near to the positive supply and its complementary Addressing Row near to the positive supply or near to ground respectively. (allowing an Addressing Row to float instead of pulling it close to ground would work, too.) When an Addressing Row is pulled to the positive supply, any of the Column Addressing Rectifiers whose anodes are connected to said Addressing Row will be forward biased and the voltage drop across those rectifiers will determine the resulting voltage on the Storage Columns connected to said forward biased Column Addressing Rectifiers. Storage Columns whose voltages are pulled up in this way are called “disabled.” The voltage level on said disabled Storage Column would be equal to the near to the positive supply voltage on the Addressing Row minus the forward voltage across said forward biased Column Addressing Rectifier. Disabling the Storage Rows by pulling their voltages down and disabling the Storage Columns by pulling their voltages up will reverse bias any Storage Rectifiers at the intersection of said disabled Storage Rows with said disabled Storage Columns. The Storage Rectifier, if any, at the intersection of the remaining enabled Storage Row with the remaining enabled Storage Column will be forward biased and the forward voltage drop of said Storage Rectifier, if any, would be centered around the voltage level of one-half of the positive supply. This would place the voltage level on the enabled Storage Row at one-half of the positive supply plus one-half of the forward voltage drop across said Storage Rectifier. Also, this would place the voltage level on the enabled Storage Column at one-half of the positive supply less one-half of the forward voltage drop across said Storage Rectifier. If no Storage Rectifier was present at the intersection of the remaining enabled Storage Row with the remaining enabled Storage Column then the voltage on that enabled Storage Row would be the positive supply and the voltage on that enabled Storage Column would be ground. (It should be noted that it does not matter if the storage rectifiers between the disabled rows and the disabled columns are forward biased because, as will be seen below, the storage bit sensing circuitry will not sense those Rows that have been shifted to the disabled voltage level.) The Storage Bit Sensing Column, M, is to be biased to a voltage level just below the positive supply minus one rectifier forward voltage drop. Recalling that a disabled Storage Row will be at a voltage level equal to a near to ground voltage plus a rectifier's forward voltage, this means that the Storage Bit Sensing Rectifiers, U, between the Storage Bit Sensing Column, M, and each of the disabled Storage Rows will be reverse biased and conduct no current; this will account for the state of all of the Storage Bit Sensing Rectifiers, U, except the one connected between the Storage Bit Sensing Column, M, and the one enabled Storage Row. If a rectifier is present at the intersection of the enabled Storage Row with the enabled Storage Column then that enabled Storage Row, then the voltage level on that enabled Storage Row would be at one-half of the positive supply plus one-half of a rectifier's forward voltage which would not be sufficient to forward bias the Storage Bit Sensing RectiFier and no current will flow to the output, OUT. On the other hand, if no rectifier is present at the intersection of the enabled Storage Row with the enabled Storage Column then that enabled Storage Row, absent any Storage Bit Sensing Rectifier, would be at the Voltage potential of the positive supply; this is sufficient to forward bias the Storage Bit Sensing Rectifier and current will flow to the output, OUT. Naturally, this assumes that the voltage level of the positive supply is sufficiently high to forward bias the various rectifiers as described. Also, it is preferable to keep the positive supply to as low a voltage as possible in order to minimize power dissipation and to minimize the reverse voltages on the Storage Bit Sensing Rectifiers, U, thereby reducing the potential current leakage through those Storage Bit Sensing Rectifiers that might tend to offset the output current. This leakage current would be fairly constant and equal to the leakage of a rectifier reverse biased by an amount equal to the difference between the bias voltage on the output and the voltage on a disabled Storage Row multiplied by the number of disabled Storage Rows and, as such, could be corrected for if necessary. For example, an opposite current of equal magnitude could be injected into the Storage Bit Sensing Column. As shown in FIG. 3 , to the left of each Storage Row, P, is an address designation shown as an “X” followed by some binary digits. The “X” signifies the upper address bit inputs (which control the Addressing Rows, Q,) and the binary digits signify the lower address Bit inputs which control the Addressing Columns, N. As is shown, for each address bit input there are two Addressing Columns—one corresponding directly to the address bit input (labeled An) and one corresponding to the complement of that address bit input (labeled −An). Each bit position in the address designation of any Storage Row corresponds to a given address bit input and to a complementary pair of Addressing Columns. When the bit position contains a 0, a rectifier is connected between that Storage Row and the Addressing column corresponding to the complement of that address bit input. when the bit position contains a 1, a rectifier is connected between that Storage Row and the Addressing Column corresponding directly to that address bit input. Also as shown in FIG. 3 , below each Storage Column, O, is an address designation shown as some binary digits followed by an “X”. The “X” signifies the lower address bit inputs (which control the addressing Rows, N) and the binary digits signify the upper address bit inputs which control the Addressing Rows, Q. As is shown, for each address bit input there are two Addressing Rows—one corresponding directly to the address bit input (labeled An) and one corresponding to the complement of that address bit input (labeled −An). Each bit position in the address designation of any Storage Column corresponds to a given address bit input and to a complementary pair of Addressing Rows. When the bit position contains a 0, a rectifier is connected between that Storage Column and the Addressing Row corresponding directly to that address bit input. When the bit position contains a 1, a rectifier is connected between that Storage Column and the Addressing Row corresponding to the complement of that address bit input. Refer now to FIG. 4 which shows identical circuitry to that shown in FIG. 3 except for the addition of complementary address input buffer-driver circuitry, V, W, X, and Y, as well as an output driver circuit, Z. Following the addressing of one complementary pair of addressing transistors controlled by address line A 2 , one will see that transistors Q 2 and Q 4 will be turned off through resistors R 10 and R 12 when A 2 is at a low logic state (a low enough voltage that the base emitter junctions of Q 2 and Q 4 are not forward biased) or floating. When transistor Q 4 is turned off, transistor Q 3 will be turned on by the positive supply through resistor R 11 . One will also see that transistors Q 2 and Q 4 will be turned on through resistors R 10 and R 12 when A 2 is at a high logic state. When Transistor Q 4 is turned on, it will pull away the current available at the base of Q 3 , dropping the voltage on the base of Q 3 and thereby turning off Q 3 . The result is that the directly addressed Addressing column, A 2 , will be pulled to near ground (approximately 0.2 v) when A 2 is low and Q 3 is turned on (saturated), or the complementary addressed Addressing Column, −A 2 , will be pulled to near ground (approximately 0.2 v) when A 2 is high and Q 2 is turned on. Following the addressing of another complementary pair of addressing transistors controlled by address line A 3 , one will see that transistors Q 11 and Q 12 will be turned off through resistors R 19 and R 20 when A 3 is at a high logic state (a high enough voltage that the base emitter junctions of Q 11 and Q 12 are not forward biased) or floating. When transistor Q 11 is turned off, transistor Q 13 will be turned on through resistor R 21 to ground. One will also see that transistors Q 11 and Q 12 will be turned on through resistors R 19 and R 20 when A 3 is at a low logic state. When transistor Q 11 is turned on, it will cause the voltage on the base of Q 13 to come within 0.2 v of the positive supply and thereby turning off Q 13 . The result is that the complementary addressed Addressing Row, −A 3 , will be pulled to within 0.2 v of the positive supply when A 3 is low and Q 12 is turned on, or the directly addressed Addressing Row, A 3 , will be pulled to within 0.2 v of the positive supply when A 3 is high and Q 13 is turned on. The output driver circuit, Z, achieves the voltage biasing of the Storage Bit Sensing Column, M, while providing gain to the output current. The three rectifiers, D 33 , D 34 , and D 35 , serve to set the voltage on the emitter of transistor Q 1 at three rectifier forward voltage drops below the positive supply. Assuming that the forward voltage drop across the base-emitter junction of transistor Q 1 is the same as the forward voltage drop of the rectifiers, this would require that the base of that transistor be at a voltage level no less than two rectifier forward voltage drops below the positive supply in order for its collector to draw current, which in turn would require that the one enabled Storage Row be at a voltage level no less than one rectifier forward voltage drop below the positive supply in order for the collector of transistor Q 1 to draw current. Naturally, it is not necessary that the forward voltage drop across the base-emitter junction of transistor Q 1 would be the same as the forward voltage drop of the rectifiers and one skilled in the art will know of other ways to bias the Storage Bit Sensing Column. The circuit used to describe the operation of a DRS Array and shown in FIG. 4 can be constructed using silicon diodes (such as the 1N914) for all of the rectifiers; 1 MΩ resistors for all of the Row Resistors, R, and Column Resistors, S, and for the resistor in the output driver circuit, R 9 ; 10KΩ resistors for all of the resistors of the address input buffer-driver circuitry, V and W; 2N3904 transistors for all of the NPN transistors and 2N3906 transistors for all of the PNP transistors. The output will sink virtually no current when addressing the point of intersection of a Storage Row with a Storage Column that has a rectifier present at that point. If no rectifier is present at the addressed point of intersection, approximately 0.7 μA will flow into the base of Q 1 and, assuming a current gain (β) of about 100 times for Q 1 , the output will therefore sink abouto 0.07 mA. Some variation on this idea will be apparent. FIG. 5 shows two variations on the Dual-addressed Rectifier Storage Array. The first variation relates to the number of bits of stored information that will be retrieved at the same time. Notice that Row Addressing Rectifiers D 9 through D 16 , transistors QZ through Q 4 and resistors R 10 through R 12 have been removed, that transistor Q 23 , rectifiers D 86 , D 87 , and D 88 and resistor R 39 have been added, and that the Storage Bit Sensing Rectifiers have been split into two groups where the cathodes of rectifiers D 5 through D 8 are now connected to the base of transistor Q 23 through a new Storage Bit Sensing Column. The resulting circuit is the equivalent of two DRS Arrays placed one above the other with each having four Storage Rows and eight Storage Columns where those eight Storage Columns (and their associated addressing means) are common to both arrays and where the Addressing Columns that are connected to the cathodes of the Row Addressing Rectifiers of the two arrays (and their associated addressing means) are also common to both. Transistor Q 1 , rectifiers D 33 , D 34 , and D 35 , and resistor R 9 now operate with Storage Bit Sensing Rectifiers, D 1 , D 2 , D 3 , and D 4 , to detect state of the addressed bit within one array and transistor Q 23 , rectifiers D 86 , D 87 , and D 88 , and resistor R 39 now operate with Storage Bit Sensing Rectifiers, D 5 , D 6 , D 7 , and D 8 , to detect state of the addressed bit within the other array. In this way, two bits—one from one array and one from the other array—will be read out at the same time. One skilled in the art will recognize that this variation can be applied multiple times such that even larger numbers of bits can be read simultaneously. Also shown in FIG. 5 is a variation relating to the power requirements of the device and, in particular, to limiting the current in the device to only a portion of that device thereby reducing the overall power consumption. Notice that Column Addressing Rectifiers D 36 through D 43 , transistors Q 17 through Q 19 and resistors R 25 through R 27 have been removed, that transistors Q 20 , Q 21 and Q 22 , and resistors R 36 , R 37 , and R 38 have been added and that the Storage Columns have been divided into two groups where one group is connected to the collector of transistor Q 20 through column Resistors R 28 through R 31 and the other group is connected to the collector of transistor Q 21 through Column Resistors R 32 through R 35 . Address line A 4 controls transistors Q 20 , Q 21 , and Q 22 such that Q 20 and Q 21 are a complementary pair that controls the enabling of the left four lines or the right four lines of the Storage columns. Since a Storage Column is otherwise disabled by applying a voltage to that line (through an Addressing Row and a Column addressing Rectifier) which results in a current equal to roughly the difference of the positive supply less 0.9 v divided by the Cathode Resistor value 1 MΩ (4.1 μA if using a 5 v supply), wherever a line can be disabled by cutting off the current through the Column Resistor, that much currert will be saved. These savings will become significant in much larger arrays having a very large number of Storage Columns. The limit on the number of address lines that will control this type of power reducing means for enabling Storage Columns relates to the complexity of the circuit by increasing the number of transistors (a complex component) while reducing the number of rectifiers (a relatively less complex component). One skilled in the art will recognize that this same type of power saving enabling means could also be created on the Storage Rows by using PNP transistors to control the cutting off of current to those Storage Rows (through the Row Resistors). Another variation might be to use differently made components. For example, in the above explanation, silicon rectifiers having a forward voltage of about 0.7 v and transistors having base-emitter forward voltages also of about 0.7 v were assumed. But, other types of rectifiers and transistors could be used having lower or higher forward voltages. In another variation, the Storage Rectifiers used could be the base-emitter junctions of NPN transistors (having all of their collectors tied together in the same way one would tie together the outputs of opened-collector gates in a logical AND configuration); this would make it possible to eliminate all of the Storage Bit Sensing Rectifiers, transistor Q 1 , resistor R 9 and rectifiers D 33 , D 34 and D 35 and instead sense the current sunk on the combined collectors. Since no more than one Storage Rectifier will be conducting at any time, current being sunk on the combined collectors would indicate that a Storage Rectifier (a base-emitter junction) exists at the point of intersection of the enabled Storage Row and the enabled Storage Column. Other variations could include, memory cells comprising a fusible link resulting in a One Time Programmable Read Only Memory (OTPROM) device (see U.S. Pat. Nos. 4,312,046 and 4,385,368), or memory cells comprising charge-storage devices to create a Random Access Memory (RAM) device (see U.S. Pat. Nos. 3,626,389 and 3,838,405). The above description of the preferred embodiment makes reference to a variation on the resistive means—that of using “leaky” diodes. FIG. 12 shows a plot depicting the voltage/current relationship of a typical diode, where the horizontal axis is voltage, V, and the vertical axis is current, I. Moving right from the plot's origin, one can see that no significant current flows through the diode until the Forward Voltage, F, is reached. Moving left from the plot's origin, one can see that only a small amount of leakage current flows until the Reverse Breakdown Voltage, R, is reached. However, there is a region, A, between the origin and the point of the Reverse Breakdown Voltage where the device behaves similarly to a resistor. FIG. 13 shows a variation wherein the resistive means comprises diodes that are believed to operate in this resistor-like region. The diode biasing circuit, A (which would typically not be constructed as a part of an integrated circuit comprising the remainder of the circuit shown in this figure), comprises circuitry for biasing a row of diodes, H, in this resistor-like region for the row lines as well as circuitry for biasing a row of diodes, J, in this resistor-like region for the column lines. All of the diodes would be formed as a part of the same integrated circuit and would therefore have operating characteristics that are matched. On the row lines, the diodes, H, will be biased between the voltage of a disabled row and the voltage present at the output of the operational amplifier, B. The operational amplifier circuitry will ensure that the voltage across feedback diode D 102 will be the same as is across the diodes, H. As a result, the current through diodes, H, and diode D 102 will be the same as a result of the device matching. This is accomplished by pulling down point K in identical fashion to when pulling down any of the row address lines (A 0 through −A 2 ). This ensures that the voltage on the anode of diode D 103 will match the voltage on any one of the disabled row lines since an operational amplifier operates (as one skilled in the art knows) such that the same voltage will be present at the ‘+’ and ‘−’ terminals in normal feedback mode. Current sources C and D have equal magnitude and ensure that the same current flows through D 102 and D 103 and, therefore, that the voltage drop across D 102 and D 103 is matched; they could be constructed as a part of a standard current mirror circuit. The result is that the voltage at the anode of diode D 102 matches the voltage on one of the disabled row lines, that the voltage on the anodes of diodes D 86 through D 93 matches the voltage at the anode of diode D 102 (they are connected to the row lines), and the cathodes of all the diodes D 86 through D 93 as well as the cathode of diode D 102 are at the same voltage (which is greater than that of the disabled row lines resulting in a reverse voltage across all the diodes). Given that the reverse voltage drop across diodes D 86 through D 93 is the same as that across diode D 102 , the current leaking through the diodes D 36 through D 93 will be the same as the current leaking through diode D 102 (and controlled by current source C) due to their matched operating characteristics. The same mechanism is essentially accomplished for the column line resistive means with via operational amplifier E, current sources F and G, diodes D 94 through D 101 , feedback diode D 104 and diode D 105 , but with opposite polarity. Note that point L would be sourced in identical fashion to the sourcing of the column address lines (−A 3 through A 5 ). FIG. 13 also shows a possible solution should the cumulative diode leakage currents in extremely large arrays of Storage Rectifiers become noticeable. An additional Addressing Row, N, has been included. Notice this row has a rectifier connecting it to every one of the Storage Columns and that asserting this row by pulling it to the positive supply will disable the one remaining Storage Column that had been left enabled. As a result, it is believed that one could detect the presence of a Storage Rectifier at the point of intersection of the enabled Storage Row and Storage Column even with the existence of such leakage currents by sampling the level at the output both while asserting and not asserting this Addressing Row, N. Since the presence of a Storage Rectifier will result in a loading of the one enabled Storage Row by connecting that row through the Storage Rectifier to the remaining enabled Storage Column, disabling that Storage Column would remove this loading and result in a slightly higher voltage at the Output. However, if no Storage Rectifier is present, the level at the Output should not show any significant change. Other variations might include reversing the polarity and types of some of the components. Because of the symmetry of the DRS array, many of the techniques shown in one area of the device can be implemented in the opposite area with only a reversal of polarities. for example, one skilled in the art will recognize that The Storage Bit Sensing Rectifiers could be connected to the Storage Columns if transistor Q 1 was of the PNP type, rectifiers D 33 , D 34 , and D 35 were reversed in polarity and connected in series to ground instead of to the positive supply, and resistor R 9 was connected to the positive supply instead of to ground; in this variation, the collector of transistor Q 1 would be a current source instead of a current sink. In another variation the means of row or column selection by disabling certain rows or columns might be used in combination with other selection means disclosed in the prior or subsequent art. It is believed that Dual-addressed Rectifier Storage Arrays will typically be fabricated as an integrated circuit; a possible layout of the Rows, the Columns, and the rectifiers are shown in FIGS. 6 through 10 (the addressing transistors and resistors, the Row and Column Resistors, and the connection pads to the chip have been omitted for clarity—the fabrication of these devices are well known to those skilled in the art). In FIG. 6 , an N-type wafer is shown to be doped with P-type channels which form the Rows. Next ( FIG. 7 ) is shown that an oxide layer is grown and then ( FIG. 8 ) openings corresponding to the rectifiers in the circuit are etched through that oxide layer. The data stored in the device is programmed during this step—wherever a Storage Rectifier is desired, a hole is etched through the oxide layer (it should be noted that the pattern of the openings matches the pattern of the rectifiers as drawn in FIG. 3 ). FIG. 9 shows the chip with an aluminum metalization layer and finally ( FIG. 10 ) shows the result of etching that aluminum into vertical lines which form the Columns. Wherever the aluminum contacts a doped region through one of the holes in the oxide layer, a rectifier (of the type sometimes referred to as a metal-on-silicon junction type or as a Schottky Diode type) is formed. The N-type wafer substrate would be kept at the most negative voltage in the circuit thereby creating reverse biased p-n junction between the doped regions and the substrate, the result of which is to electrically isolate those generally parallel doped regions. The generally parallel metalized regions which form the Columns are electrically isolated from each other due to their being formed upon the non-conducting oxide layer (except where they contact the doped regions through the holes in the oxide layer). Where the aluminum Columns contact the P-type Rows, a metal-on-silicon junction rectifier is formed such that the current flow (where conventional current flow is the flow of holes, that is to say, current flowing from positive voltage potential to negative voltage potential) is from the doped regions to the metalized regions when the junction is forward biased. As shown in FIG. 10 , space is available for the manufacture of the addressing components and the Row and Column Resistive means at the lower left corner of the chip and around the edges. The addressing components might also be modified when constructing a DRS Array as an integrated circuit. The pair of resistors directly connected between each address input and the bases of the two addressing transistors could be replaced by a single resistor (as shown in FIG. 3 , a single resistor could do the job of resistors R 10 and R 12 , a single resistor could do the job of resistors R 13 and R 15 , and so on). Those transistors could be external to the integrated circuit form of the DRS Array. A variation on the semiconductor manufacturing of a DRS Array might be to dope N-type regions into a P-type wafer thereby reversing the polarity of the metal-on-silicon junction rectifiers (the Rows and the Columns would be reversed). Another variation would spread out the Addressing Columns across the width of the chip—alternating Addressing Columns with one or more Storage Columns—instead of grouping those Addressing Columns at the left side of the chip; this will spread out those Lines carrying most of the current in the circuit thereby more evenly distributing the power dissipation across the chip (the same technique could be done with the Rows). Another variation would be to construct the device with p-n junction rectifiers or a combination on metal-on-silicon junction rectifiers and p-n junction rectifiers. Another variation might enable programming the stored data during the metalization etching step. Referring to FIG. 11 , an opening has been etched through the oxide layer at every potential Storage Rectifier location. Programming of the stored data bits is accomplished when the metalization layer is etched. In those locations where a storage Rectifier is desired, a metal connection is left during the metal etching step between the metal pad covering the opening in the oxide layer and the metal Column; where no storage Rectifier is desired, that metal connection is etched away. It is believed that this approach will enable all of the semiconductor manufacturing steps, except the final metal etching step to be performed and that wafers so made could be stored safely under the protective metalized layer. In this way, wafers could be mass produced without regard to the data to be stored in the chip. Economic concerns may drive several other possible variations on the semiconductor manufacturing of a DRS Array. It is expected that a significant part of the cost to manufacture these devices will be in the packaging where the greater the number of electrical connections between the package and the controlling device, the greater the cost of that packaging; a package with more leads will be more expensive and more prone to mechanical failures. As a result, it is believed that steps will be taken to reduce the number of package leads needed. One skilled in the art will quickly realize the circuitry needed to implement any of the following variations. One possibility would be to address the lower address lines directly but to retain the higher address lines internally; in other words, the latch shown in FIG. 2 would be incorporated into the DRS Array integrated circuit. Another possibility would be to incorporate a shift-register where the address bits would be shifted into the device serially and retained thereby reducing the number of addressing leads to two (the shift-register serial input and the shifting-clock). However, a more practical variation, shown in FIG. 14 , might be to incorporate both a counter and a shift register where the shift register would be used to enter an address onto the chip which could then be loaded into the counter. In this last variation of the serially loaded counter, the address would be retained within the counter so the many address lines are essentially reduced to four: one for the serial address input (S), one for clocking the shift register (K), one for clocking the counter (C), and one for loading the address into the counter from the shift register (L). One skilled in the art will quickly realize the circuitry needed to implement any of these variations. One well skilled in the art will recognize that a reduction to three lines can be achieved here resulting in a total of six connections to the chip (three plus power, ground, and data out). Finally, a variation on the digital logic device comprising one or more DRS Arrays of which one or more may be removable which themselves comprise the above mentioned serially loaded counter logic. By limiting the manufacture of such devices to having output in analog format only, the risks to the makers of programming (e.g., music and video programs) will be reduced. With CD-ROM technology, the output from some CD-ROM readers is in a digital format. As a result, any copy made will be of the same quality as the original. This potentially results in significant lost revenue as the users of this technology could casually make copies for friends and relatives that cannot be distinguished from the originals (this was not the case with prior technologies such as cassette tapes and video tapes where each successive copy degraded somewhat). By limiting the manufacture of devices comprising DRS Arrays that are addressed via serially loaded counters to analog output only, the same degradation of copies will occur thereby reducing some of the risks to the makers of programming by causing the copies to be less desirable than the originals. While devices comprising DRS Arrays that are addressed via serially loaded counters could be limited to analog output only, they could still include means for reading DRS Arrays in other formats (i.e., DRS Arrays directly addressed via many address lines), however, devices comprising DRS Arrays that are directly addressed via many address lines and which give digital access to the information stored therein would not include means for reading DRS Arrays that are addressed via serially loaded counters. It is believed that minor flaws in the semiconductor wafer will mostly impact the operation of the metal-on-silicon junction rectifiers when they are reverse biased by lowering the reverse breakdown voltage. Since the DRS Arrays are expected to be operated at low voltage levels, large reverse voltages are considered unlikely in normal operation. As a result, the impact of these minor flaws are not expected to impact the operation of the device and high device manufacturing yields are anticipated. However, the addressing transistors will likely be adversely affected by such flaws. There will likely be a trade-off between the increased cost of manufacture resulting from lowered device yields (as a result of the impact of semiconductor flaws on the increased complexity of the addressing circuit) and the savings on packaging (as a result of reducing the number of device leads). The selection of a line by disabling all undesired lines could be utilized in many related electronic devices. The rectifiers at the storage locations could be fabricated as Light Emitting Diodes LED's) with such a rectifier present at every storage location. In this way, the device could be used as a display panel where a given display pixel could be turned on by selecting that bit location; the display panel would be scanned, selecting and illuminating bit locations in sequence while skipping bit locations that are to remain dark. Also, using a technique such as pulse width modulation, which is well known to one skilled in the art, one could even control the duration of a pulse of light emitted at any given pixel location and thereby control the perception of the intensity of the light emitted. The high expected storage densities come from the symmetry the design—the Storage Bit Sensing Rectifiers, The Addressing Rectifiers, and the Storage Rectifiers are all constructed in the same way. The result of this is that they can all be made at the same time with the same semiconductor manufacturing steps. By using metal-on-silicon junction rectifiers, the primary components in the circuit are essentially constructed vertically on the semiconductor's surface instead of horizontally as might conventionally be done resulting in a very efficient use of the semiconductor “real estate”. The scaling up of the device is expected to be easily accomplished. For example, on a one inch square chip, if the Anode Lines and the Cathode Lines are placed at roughly 0.45 micron center to center spacing, then a DRS Array that is roughly 65,536 by 65,536 Lines could be made. This is the equivalent of about 4,294,967,296 bits or about 536,870,912 bytes or about the capacity of a present day CD ROM. State of the art technology at the time of this writing is below 1 micron line widths. It is envisioned that the present invention could be used anywhere that one would use a CD-ROM drive or player or anywhere a large amount of information is needed. The foregoing description of an example of the preferred embodiment of the invention and the variations thereon have been presented for the purposes of illustration and description. It is not intended to be exhaustive or to limit the invention to the precise forms disclosed. Many modifications and variations are possible in light of the above teaching. It is intended that the scope of the invention be limited not by this detailed description, but rather by the claims appended hereto.

A read-only data storage and retrieval device is presented having no moving parts and requiring very low power. Addressing can be accomplished sequentially where the address increments automatically or can be accomplished randomly. High density storage is achieved through the use of a highly symmetric diode matrix that is addressed in both coordinate directions; its symmetry makes the Dual-addressed Rectifier Storage (DRS) Array very scaleable, particularly when made as an integrated circuit. For even greater storage flexibility, multiple digital rectifier storage arrays can be incorporated into the device, one or more of which can be made removable and interchangeable.

BACKGROUND OF THE INVENTION [0001] 1. Field of the Invention [0002] This invention relates to communication networks, and particularly to a network that provides portable users with secure access when exchanging information with other users on the network. [0003] 2. Discussion of the Known Art [0004] As military conflicts are being resolved more through the use of a network-centric rather than a platform-centric paradigm, vital communications over the established networks must be secure, reliable, interoperative, survivable, and timely. The implementation of a high capacity, multimedia network is also desirable. [0005] Free space optical (FSO) or photonic communication links have been deployed in fixed, point-to-point links for commercial and military applications. Such links may be preferred over microwave or millimeter wavelength radio frequency (RF) links for short range communications, especially when other communication infrastructure is unavailable, unreliable, or untrustworthy. FSO links have the following advantages: [0006] 1. The links are highly directional, and therefore quite immune to interception, interference or jamming. [0007] 2. Secure communications during periods of radio silence. [0008] 3. Elimination of any detectable RF signature. [0009] 4. FSO terminals can be made small, lightweight, and are easily portable. Optical antennas including light emitters (e.g., laser LEDs) and detectors (e.g., photodiodes) have typical gains on the order of one million times those of isotropic RF antennas. [0010] 5. Low power consumption. [0011] 6. The availability of a wide frequency spectrum with no governmental regulatory restrictions. [0012] 7. Large data bandwidth capacity. [0013] 8. Direct baseband signaling, thus simplifying modulation and demodulation processes. [0014] 9. Ease of multiplexing, de-multiplexing, and switching of optical channels. [0015] 10. Tactically useful range. [0016] Projects are being pursued that would enable laser communication on the move between platforms ground to ground, ground to air, air to air, air to satellite, and satellite to air. Infrared (IR) light sources and detectors suitable for use in high data rate FSO transmitters and receivers are commercially available at low cost. [0017] IR light penetrates clear glass but will not propagate through walls or other opaque building structures. FSO links are therefore confined to rooms or other areas inside buildings where the links are established. Such confinement enhances the security of FSO transmissions against interception or casual eavesdropping, and avoids interference between optical links operating in physically separate regions, thus making possible a high degree of spectrum reuse. Also, while multipath fading may cause signals to fluctuate in strength and phase over RF links, FSO links are immune to fading if intensity modulation and direct detection (IM/DD) techniques are applied. See, J. M. Kahn et al., “Wireless Infrared Communications”, 85 (2) Proceedings of the IEEE (Feb. 1997), at 265-98, which is incorporated by reference. Portable Infrared (IR) Devices [0018] For short range (up to a few meters) applications, consumer devices are available that allow data to be transferred between the devices via infrared light. The Infrared Data Association (IrDA) defines specifications for point-to-point communication using directional half duplex serial IR links through space, at data rates up to and including 115.2 kbit/s; 0.576 Mbps, 1.152 Mbps, 4.0 Mbps and 16 Mbps. Cell phones are available with IR ports that follow these standards for enabling the phones to dump data into stationary printers, PDAs, or PCs equipped with IR ports. See “Motorola i930/i920”, at <www.phonescoop.com/phones/phone.php?p=627>. IR ports of typical cell phones do not carry active voice communications and, as mentioned, are limited in range to 1 to 2 meters. VoIP Telephony and Wireless Local Area Networks [0019] The use of voice-over-Internet protocol (VoIP) telephony, both wired and RF wireless, is expanding. In a conventional circuit switched telephone system, a dedicated physical connection is established between a calling and a called party over the duration of the call. The continuous connection assures that voice signals carried between end points of the system are not interrupted. With a VoIP system, however, there is no dedicated connection. Instead, analog voice signals from a microphone transducer in a user's handset or headset are digitized, and corresponding digital data is transmitted over a system network in separate groups of data called “packets”. Each packet contains the sender's and the recipient's IP addresses, and a piece of digitized voice information (“payload” data). The packets may be routed through the network over different paths, and eventually arrive with some delay at a common destination to be recombined in the proper sequence. Further, each packet may arrive with a different delay. Variations in arrival time are defined as “jitter”. Some packets may never reach the destination, resulting in “packet loss”. Most vendors adhere to strict limits on tolerable packet loss, delay, and jitter. For example, Cisco Systems adopted the following guidelines for VoIP network operation: [0000] Network Performance Value Delay <=150 milliseconds (ms) one-way Jitter <=30 ms Packet loss <=1% [0020] VoIP may offer many features above and beyond those afforded by traditional telephony systems, whether wired or remote. See, e.g., A. Noser, “Combining VoIP and Wireless Services”, at <www.ncstate.net/wireless/presentations/wirelessvoip/wirelessvoip.html>, which is incorporated by reference. Manufacturers claim their wireless VoIP products allow mobile users to engage in conversations anywhere in an IP network with reliability and voice quality equivalent to that of a desktop office phone. Internet gateways and RF access points are positioned to ensure that user conversations do not drop out or experience gaps, regardless of a user's location within a defined area. As voice quality, reliability and security improve, IP wireless communication including the use of convenient portable VoIP handsets is likely to increase. [0021] A typical VoIP local area network (LAN) 10 is illustrated in FIG. 1 including commercial off-the-shelf (COTS) products. To connect with a legacy public telephone switching exchange (PBX) 12 , a telephony gateway 14 is configured to convert analog voice signals received over the PBX 12 into IP voice data packets. The packets are routed through an Ethernet cable 15 that connects with RF wireless access points 16 . Voice data packets arriving at the gateway 14 over the cable 15 are converted to analog voice signals for transmission into the PBX 12 . The gateway 14 may be omitted if the PBX 12 is a so-called telephony server. [0022] The access points 16 may comprise RF wireless routers each of which operates according to, e.g., known IEEE 802.11x signaling protocols. A voice priority server 20 available, for example, from Spectralink SVP® may be provided to ensure that the voice data packets have priority over other kinds of data carried over the network 10 . The access points 16 may join or bridge various wireless clients such as, for example, a number of portable VoIP telephone sets 26 , a notebook computer 27 and a PDA 28 , with fixed users and devices connected by wire to the network 10 . [0023] FIG. 2 shows a typical high level architecture for a wireless access point 16 . Access point 16 may operate, for example, under one or more defined RF signaling protocols per the IEEE Standards 802.xxx. Because voice data transmitted by a user of a RF device may be received by users of like RF devices within range, some security measures are available to ensure that a user's data is not captured or manipulated by unauthorized intruders. When classified or other highly sensitive voice messages are involved, however, commercial security (COMSEC) is insufficient for the task. For example, adding improved Type I security can significantly increase cost and management complexity, since such security must be controlled and crypto keys must be managed. Wireless VoIP Phone Sets [0024] Several vendors provide RF wireless VoIP telephone sets that can access a LAN using IEEE 802.11x or other newly emerging IEEE 802.xxx RF signaling protocols. See, e.g., “Linksys WIP330 Wireless-G IP Phone”, at <www.tomsnetworking.com/Sections-article147-page1.php>. A block diagram of a typical wireless VoIP telephone 26 is shown in FIG. 3 . Core subsystems include: [0025] A RF transceiver/power amplifier 30 that performs frequency translation between the RF and the baseband (voice) signals, and amplifies RF signals to be radiated from the phone from an antenna 31 . [0026] A medium access control (MAC)/baseband processor 32 which implements the applicable IEEE 802.xxx protocols and provides modem functionality to control wireless signaling and communication between the telephone 26 and the wireless access points of the LAN. [0027] A DSP/microcontroller/OMAP 34 that executes VoIP call controls and voice processing, and provides a user interface. [0028] Various memories including flash, ROM and RAM stages for storing programming code, voice and other data. [0029] A voice coder-decoder (CODEC) 36 which interfaces with a user headset 37 having a microphone 38 and a speaker or earpiece 39 . The CODEC 36 operates to convert a user's analog voice signals as produced by the microphone 38 , into corresponding digital voice data to be processed by the OMAP 34 . [0030] The RF bandwidth required for each voice call depends on (i) the type of CODEC 36 , (ii) the number of CODEC samples per data packet, and (iii) the packet header compression. The number of CODEC samples per packet affects the delay of a VoIP call. As the size of the sample data increases, the required bandwidth decreases but the overall delay increases. [0031] As mentioned, if a wireless VoIP telephone set user desires to discuss classified subject matter, COMSEC items must be provided thereby increasing equipment cost and management complexity. Accordingly, there is a need for a robust multi-user local area wireless network that is not only capable of interfacing with current VoIP telephone sets, but which also provides security for portable users who want to convey sensitive information without having to invoke costly COMSEC measures. SUMMARY OF THE INVENTION [0032] The inventive network allows a portable user to engage in wireless communications wherein normal messaging is routed over a RF link with the user, and classified or other highly sensitive messages are contained over a secure FSO link that can be established by or with the user when desired. [0033] According to the invention, a communication network includes a local area network (LAN) and a wireless access point coupled to the LAN. In one embodiment, each access point includes a medium access control (MAC) stage, and a radio frequency (RF) transmitter/receiver for communicating unsecure message data via RF links with users of associated wireless devices. An optical transmitter/receiver in the access point enables the users to communicate secure message data over the LAN via free space optical (FSO) links with the users. [0034] The MAC stage operates (i) to direct unsecure data from the LAN to the wireless device users and to direct unsecure data from the users to the LAN, via the RF transmitter/receiver; and (ii) to direct secure data from the LAN to the wireless device users and to direct secure data from the users to the LAN, via the optical transmitter/receiver. [0035] According to another aspect of the invention, a wireless handset includes a message data source, and a radio frequency (RF) transceiver for transmitting RF signals corresponding to unsecure message data to a network access point, and for receiving RF signals corresponding to unsecure message data radiated from the access point. An optical transceiver operates to transmit free space optical (FSO) signals corresponding to secure message data to an optical access antenna system associated with the access point, and to receive FSO signals corresponding to secure message data emitted from the optical antenna system. A switching stage has a first port coupled to the message data source, a second port coupled to the RF transceiver, and a third port coupled to the optical transceiver. The switching stage is configured to couple the message data source to the RF transceiver for unsecure message data, and to the optical transceiver for secure message data. [0036] For a better understanding of the invention, reference is made to the following description taken in conjunction with the accompanying drawing and the appended claims. BRIEF DESCRIPTION OF THE DRAWING [0037] In the drawing: [0038] FIG. 1 is a block diagram of a typical local area network (LAN) with wireless access points; [0039] FIG. 2 is a block diagram of a typical wireless access point in the network of FIG. 1 ; [0040] FIG. 3 is a block diagram of a typical voice over Internet protocol (VoIP) wireless telephone; [0041] FIG. 4 is a block diagram of a communication network according to the invention; [0042] FIG. 5 is a block diagram of an integrated radio frequency (RF) and free space optical (FSO) wireless access point, according to the invention; [0043] FIG. 6 is a block diagram of an integrated RF and FSO wireless handset, according to the invention; [0044] FIG. 7( a ) illustrates a first embodiment of an optical access antenna system, including a number of optical antennas associated with the access point of FIG. 5 ; [0045] FIG. 7( b ) illustrates a second embodiment of the optical access antenna system; [0046] FIG. 8 shows optical transmitting and receiving elements mounted on the handset of FIG. 6 ; and [0047] FIGS. 9( a ) to 9 ( d ) illustrate arrays of light receiving elements that may form part of each optical antenna in the system of FIG. 7( a ) or FIG. 7( b ). DETAILED DESCRIPTION OF THE INVENTION [0048] FIG. 4 is a schematic block diagram of a communication network 40 according to the invention. The network 40 has one or more associated wireless access points (WAP) 50 , described below, which enable the network to be accessed by users of one or more portable handsets or headsets 70 . In addition to signaling via RF links with the access points 50 , the handsets 70 are capable of establishing FSO links when necessary to exchange secure (e.g., classified) voice data over the network 40 . Details of the handsets 70 are set out below in connection with FIG. 6 . [0049] FIG. 5 is a schematic block diagram of a first embodiment of an integrated RF and FSO network access point 50 , according to the invention. In addition to the RF components of the typical wireless access point 16 in FIG. 2 , the inventive access point 50 includes an optical transceiver. The optical transceiver comprises an optical receiver 52 , an optical transmitter 54 , and an optical access antenna system 56 that is coupled to an input of the receiver 52 and to an output of the transmitter 54 . [0050] The optical access antenna system 56 may be coupled to the optical receiver 52 and the optical transmitter 54 through a passive optical network (PON) 57 , as shown in FIG. 7( a ). The individual optical antennas 100 may be mounted, for example, in a grid array on the ceiling of one or more secure rooms access to which is restricted to authorized personnel. [0051] In the access point 50 , a baseband output of the optical receiver 52 is coupled to an input of a medium access controller (MAC) 58 through a desired crypto device 60 . The crypto device 60 operates to encrypt voice data detected by the optical receiver 52 , and to supply the encrypted voice data to the MAC 58 . Further, the optical transmitter 54 has an input coupled to a baseband output of the MAC 58 through a corresponding decrypto processor 62 . The decrypto processor 62 is configured to decode encrypted voice data received over the LAN 10 and output from the MAC 58 , and to supply the decoded data to the optical transmitter 54 . [0052] Depending on the nature of voice data originating from the LAN 10 and destined to a particular handset user, the MAC 58 routes the data through only one of the optical transceiver ( 52 , 54 ), or the RF transceiver 64 . For encrypted secure data to be delivered from the LAN 10 to an authorized handset user, the decrypto baseband processor 62 decrypts the data before it is modulated onto a light signal by the optical transmitter 54 . Voice data originating from a handset user over his/her established FSO link, is detected by the optical receiver 52 and input to the crypto device 60 , as shown in FIG. 5 . [0053] FIG. 6 is a schematic representation of an integrated wireless RF and FSO handset (or headset) 70 , according to the invention. In addition to the components of the wireless VoIP telephone 26 in FIG. 3 , the handset 70 includes an optical transceiver 72 , an optical antenna 74 which is coupled to the transceiver 72 , and a switching stage 76 . In the illustrated embodiment, the switching stage 76 has a first port 78 coupled to the MAC/baseband processor 32 of the handset 70 , a second port 80 coupled to the RF transceiver 30 , and a third port 82 coupled to the optical transceiver 72 . The handset 70 may also feature a ringer unit 84 that is coupled to an output of the optical transceiver 72 . The ringer unit 84 is constructed and arranged to produce, for example, a distinct alert sound and a blinking red LED display when the transceiver 72 detects a light signal having message data that is addressed to a user of the handset 70 . [0054] The optical antenna 74 and the transceiver 72 may be housed together in an optical module 86 that is constructed and configured to connect with the MAC/baseband processor 32 inside the handset 70 via, for example, an RJ-45 or other common wire connector interface that has been mounted onto the handset housing. The optical module 86 may be powered, e.g., by an existing voltage source (not shown) disposed in the handset 70 . If desired, the switching stage 76 and the ringer 84 may be mounted and arranged inside an existing VoIP handset. [0055] When a user of the handset 70 wants to communicate classified or other sensitive information to an authorized person on the network 40 , the user operates the switching stage 76 to establish a FSO link between the handset antenna 74 and one or more of the optical antennas 100 in line of sight of the user. As mentioned, the FSO link provides communication security since the user's light signals will not propagate beyond the room or other area in which the user and the optical antennas 100 are located. A password may be entered by the user before the switching stage 76 can be operated to establish the FSO link. It is also preferable to configure the switching stage 76 so that only one of an RF or an FSO link can be established by the user at any given time. Thus, once an FSO link has been selected, there is no possibility of an inadvertent leakage of the user's secure information onto an RF link with one of the network access points 50 . [0056] As mentioned, the optical access antenna system 56 may include a grid of the individual optical antennas 100 mounted, for example, on the ceiling of a restricted occupancy room or other limited access area in a building. In the embodiment of FIG. 7( a ), the passive optical network (PON) 57 may include one or more large core (e.g., >100μ) multimode optical fibers to couple the optical transceiver in the access point 50 with each of the optical antennas 100 forming the grid. Light reflectors or diffusers (not shown) may, if necessary, be provided in a given room to obtain 100% FSO connectivity for authorized users at various locations in the room. A variety of common building materials may also act as efficient diffuse infrared reflectors. For example, in the 800 to 900 nm range, plaster walls and acoustical ceiling tiles have diffuse reflectivities typically in a range between 0.6 and 0.9. [0057] The PON 57 in FIG. 7( a ) may, for example, implement a known coarse wavelength division multiplexing (CWDM) scheme. The CWDM scheme maintains large spectrum separation between the transmitting and the receiving light signals, so that available optical filters with high isolation can be used to separate the transmitted and the received light signals from one another at both ends. For example, to permit the use of low cost, large area silicon diode based detectors, a high power 950 nm laser may be used as a light source at the access point transmitter 54 for downstream (access point to user) transmissions, and an 880 nm GaAs laser/LED may be used as a light source for the transceiver 72 in the handsets 70 for upstream transmissions. The 950 nm laser can deliver up to 1 watt (W) of power which is sufficient to feed the multiple optical antennas 100 , thus eliminating the need for an optical amplifier. Because eye safety is of paramount importance, however, a 950 nm wavelength may not be suitable for all applications. In such case, a 1550 nm laser may be used together with an optical amplifier to increase power level. Alternatively, an element such as a diffuser may be employed to destroy special coherence of the laser beam and spread the radiation over a sufficiently extended aperture and angle. [0058] A second embodiment of the optical access antenna system 56 is shown in FIG. 7( b ). In the embodiment of FIG. 7( b ), the access point 50 is preferably located in the same room or other restricted area as the array of optical antennas 100 forming the optical access antenna system 56 . An electrical wire or cable distribution system 67 is arranged to couple the input of the crypto device 60 and the output of the decrypto baseband processor 62 in the access point 50 , with pairs of electrical to optical (E/O) media converters 102 , 104 . Each pair of E/O converters is associated with a given one of the antennas 100 . The E/O converter 102 is configured to convert electrical signals from the decrypto baseband processor 62 , into corresponding light signals to be emitted from the associated optical antenna 100 on an FSO link. The E/O converter 104 is configured to convert light signals received by the antenna 100 on the FSO link, into corresponding electrical signals for input to the crypto device 60 . The E/O converters 102 , 104 may incorporate suitable LEDs in the 880 to 1550 nm wavelength range for the uplink (E/O converter 104 ) and the downlink (E/O converter 102 ) message data flows. Because decrypted electrical data signals may be present on the cable distribution system 67 in FIG. 7( b ), it is important that appropriate measures are taken to prevent unauthorized access or detection of any signals on the distribution system 67 . [0059] If the handsets 70 include the mentioned type G.729 codecs with compressed data headers and 100 simultaneous system users are assumed, less than 1.2 Mbit/s total bandwidth is needed in each direction for voice traffic. The light source in each antenna 100 may then take the form of a LED, a Fabry Perot (FP) broad area laser, or a GaAs VCSEL based transmitter, all of which can support the mentioned data rate. [0060] In FIGS. 7( a ) and 7 ( b ), the optical antennas 100 are arrayed so as to enable a handset user to have a clear LOS to at least one of the antennas from any location in a given secure area. Because the PON architecture of FIG. 7( a ) requires no active optical components between the access point 50 and the optical antennas 100 , micro-cells 102 each of radius R less than, e.g., ten feet, may be defined. All cells may be in the same building, or spread over different buildings/rooms. For example, if the size of the building room in FIG. 7( a ) is about 80′×40′, it may be divided into 15 micro cells each with a radius R of eight feet. If a downstream laser from the access point transmitter 54 produces 100 mW of power and is split through a series of 1×4 splitters as shown in the drawing, at least 5 mW of power will be available at each optical antenna 100 . The available antenna power may then be split further to feed four or five light transmitting elements that define each antenna 100 for covering all directions. Each antenna element will then radiate about 1 mW power for downstream optical signals after discounting any losses in the PON 57 . About 1 mW of power may also be satisfactory for upstream optical signals transmitted from an antenna element on the handset 70 (see FIG. 8) . [0061] In the arrangement of FIG. 7( b ), LEDs can be used as transmitting elements for both uplink and downlink, with each LED emitting more than 1 mw power. The antenna grid in FIG. 7( b ) will not, however, be passive since the pairs of E/O converters 102 , 104 associated with each antenna 100 will require electrical power supplied, e.g., from the access point 50 in order to operate. [0062] In some applications it may also be desirable to employ an optical concentrator or lens to increase the effective area of each optical antenna 100 . An angle-diversity receiving array using multiple receiving elements 120 oriented in different directions together with a light concentrator, may be used advantageously in place of a single receiving element as shown in FIG. 9( a ). This scheme allows the receiving elements 120 to achieve high optical gain and a wide field of view (FOV) simultaneously, and may also reduce the impact of any ambient light noise and multi-path distortion. Multiple signals may be summed with equal weights, or the signal having the best signal to noise ratio (SNR) may be selected by operation of a selector/combiner stage 122 . [0063] FIG. 9( b ) shows an alternative arrangement to implement angle-diversity reception, using an array of photo detector elements 130 disposed at a focal plane of an optical concentrator 132 . Each detector element has an associated preamplifier 134 , and the elements 130 can be fabricated in large number monolithically. Only one concentrator 132 may be needed regardless of the number of detector elements 130 . The FIG. 9( b ) arrangement results in a narrower FOV as shown in FIG. 9( d ), when compared to the FOV in FIG. 9( c ) obtained when using the receiving elements 120 in FIG. 9( a ). [0064] For upstream light signals to be beamed from the handsets 70 to one or more of the optical antennas 100 , any of the mentioned devices capable of emitting light at wavelengths of 850 nm to 1550 nm may be used for the handset transmitting element 112 . Typical packaged LEDs emit light into semi-angles (at half power) ranging from about 10 to 30 degrees, making them suitable for directed transmissions. A disadvantage of LEDs is their broad spectral width (typically 25 to 100 nm) which would require a wide passband for the light detectors that define the optical antennas 100 in FIG. 7( a ), resulting in poor rejection of the ambient light. An array of available, low cost 850 nm VCSELs may therefore be useful to form directive light beams to carry the upstream signals from the handsets 70 in place of the single transmitting element 112 . For ease of implementation and to prevent inter symbol interference due to different times of arrival of voice data from a handset user, it may be desirable to use short pulse (RZ type) on-off key modulation, NRZ, or 4-PPM. [0065] Experimental results reported in the literature suggest that the above mentioned power levels for the light sources in the access point 50 and the handset 70 , will provide adequate margins to support a data rate of about 5 Mbps using a 10 mm aperture for the handset receiving elements 110 in FIG. 8 . Ultimate system performance will, of course, be limited by ambient noise and noise suppression methods. [0066] Intense ambient IR noise in the environment of a handset user may be reduced through optical filtering and/or the use of a directional light receiving array on the handset 70 to discern a desired signal from the noise. FIG. 8 shows a quadrant array of light receiving elements 110 for high collection efficiency, and a central light transmitting element 112 . The elements 110 , 112 may be mounted together, for example, on an outside surface of the handset housing or on an associated headset. [0067] Multi-megabit capacity FSO links may therefore be established by portable users on the network 40 , and known time-division multiple-access (TDMA) techniques may be applied to share available bandwidth so that a number of independent voice streams will be supported simultaneously. Some level of security may also be obtained for RF links carrying unclassified voice communications between the handsets 70 and the access point 50 , by using VoIP phones that incorporate known secure socket layer (SSL) technology. As mentioned, the switching stage 76 is preferably configured so as to make it impossible for the handset 70 to establish an RF link once an FSO link has been selected for secure communication. [0068] It will be understood that final configurations of the handset optical antenna 74 , and the optical access antenna system 56 , will depend on the physical size and nature of the building in which the antenna system 56 is installed and the number of handset users, among other parameters. Because the voice data is preferably IP in nature and the FSO links allow a large data carrying capacity, the same architecture will support multimedia services (voice, image, and other kinds of data) seamlessly, if needed. [0069] The inventive communication network 40 integrates optical communication techniques with emerging commercial VoIP handset technology. The network features secure photonic voice links including, if desired, a TDMA access scheme for classified audio transport within restricted areas. The network may therefore support any service (voice, data or image) now supported by existing RF wireless VoIP phone sets. [0070] While the foregoing description represents preferred embodiments of the invention, it will be obvious to those skilled in the art that various changes and modifications may be made without departing from the spirit and scope of the invention as defined by the following claims. For example, the network 40 may extend and enhance any existing military (e.g., JTRS) or homeland security infrastructure for which a secure access feature is desired for portable or mobile users. Also, the PON 57 in the embodiment of FIG. 7( a ) may implement optical wavelength division multiplexing using two wavelengths in each direction, one wavelength being used for classified and the other for unclassified voice signals.

A communication network includes a local area network (LAN) and a wireless access point coupled to the LAN. In one embodiment, each access point includes a medium access control (MAC) stage, and a radio frequency (RF) transmitter/receiver for communicating unsecure message data via RF links with users of associated wireless devices. An optical transmitter/receiver in the access point enables the users to communicate secure message data over the LAN via free space optical (FSO) links with the users. The MAC stage operates (i) to direct unsecure data from the LAN to the wireless device users and to direct unsecure data from the users to the LAN, via the RF transmitter/receiver; and (ii) to direct secure data from the LAN to the wireless device users and to direct secure data from the users to the LAN, via the optical transmitter/receiver. An integrated VoIP/FSO portable handset is also disclosed.

CROSS REFERENCE TO RELATED APPLICATIONS This application is a divisional of application Ser. No. 11/410,098, filed Apr. 25, 2006, now U.S. Pat. No. 7,981,876 which claims priority to Provisional Application 60/674,531, filed Apr. 25, 2005, which are hereby incorporated by reference in their entireties. This application is related to U.S. patent application Ser. No. 10/268,660 (filed Oct. 11, 2002) and to U.S. Pat. No. 5,192,756 (issued Mar. 9, 1993), U.S. Pat. No. 6,262,283 (issued Jul. 17, 2001) and U.S. Pat. No. 6,610,866 (issued Aug. 26, 2003), each of which is incorporated in its entirety by reference. FIELD OF THE INVENTION This application is directed to select squalamine salts, methods of their synthesis, their therapeutic use and their advantages relating to manufacturing, product stability and toxicity. More specifically, this application is directed to various forms of the dilactate salt of squalamine and their utility in inhibiting neovascularization and endothelial cell proliferation. BACKGROUND OF THE INVENTION Several aminosterol compositions have been isolated from the liver of the dogfish shark, Squalus acanthias . One such aminosterol is squalamine (3β-(N-[3-aminopropyl]-1,4-butanediamine)-7α, 24R-dihydroxy-5α-cholestane-24-sulfate), the chemical structure of which is shown in FIG. 1 . This aminosterol, which includes a sulfate group at the C-24 position, is the subject of U.S. Pat. No. 5,192,756 to Zasloff et al., which describes squalamine's antibiotic properties. Since its discovery, however, several other interesting properties of squalamine have been revealed. Most notably, as described in U.S. Pat. No. 5,792,635 (issued Aug. 11, 1998) and U.S. Pat. No. 5,721,226 (issued Feb. 24, 1998), which are incorporated in their entirety by reference, squalamine may inhibit the growth of endothelial cells and therefore function as an antiangiogenic agent. The use of squalamine as an antiangiogenic agent for the treatment of neovascularization in the eye and for the treatment of cancers is disclosed in U.S. patent application Ser. No. 09/985,417 (filed Nov. 24, 1998) and U.S. Pat. No. 6,147,060 (issued Nov. 14, 2000) and U.S. Pat. No. 6,596,712 (issued Jul. 22, 2003) which are also incorporated in their entirety by reference. Methods for synthesizing squalamine have been described in, for example, U.S. Pat. No. 6,262,283 (issued Jul. 17, 2001), U.S. Pat. No. 6,610,866 (issued Aug. 26, 2003), U.S. Pat. No. 5,792,635 (issued Aug. 11, 1998) and in U.S. patent application Ser. No. 10/268,660. These U.S. patents and patent applications are incorporated in their entirety by reference. Although squalamine has been previously reported to inhibit the proliferation of endothelial cells and therefore found to be useful as an angiogenesis inhibitor, a need still exists for forms of squalamine that can be readily administered to patients, especially in the form of therapeutically active, soluble salts that exhibit thermal stability upon storage and minimal toxicity and for economical methods for the manufacture of these salts. Accordingly, the identification of salts of squalamine which satisfy these requirements and which specifically inhibit angiogenesis, is an object of this invention. SUMMARY OF THE INVENTION The present invention relates to various salt forms of squalamine that inhibit endothelial cell proliferation and therefore regulate and/or modulate angiogenesis. The invention also relates to compositions which contain these salts, and methods of their use to treat angiogenesis-dependent diseases and conditions, such as, for example, cancer, tumor growth, atherosclerosis, age related macular degeneration, diabetic retinopathy, retinal ischemia, macular edema and inflammatory diseases in mammals, particularly humans. An aspect of the invention is an amorphous form or crystalline form of the dilactate salt of squalamine (3β-(N-[3-aminopropyl]-1,4-butanediamine)-7α,24R-dihydroxy-5α-cholestane-24-sulfate). In an embodiment of the invention, the crystalline form of the dilactate salt exists as a solvate. In another embodiment the crystalline form exists as a hydrate and in a further embodiment the dilactate salt exists as a solvate and a hydrate. Another aspect of the invention is a method of treating or preventing cancer in a mammal in need of such treatment, comprising administering to said mammal a therapeutically effective amount of the amorphous or crystalline forms of the dilactate salt. Another aspect of the invention is a method of treating or preventing neovascularization in a mammal in need of such treatment, comprising administering to said mammal a therapeutically effective amount of the amorphous or crystalline forms of the dilactate salt. In select embodiments, the neovascularization is in the eye, in the gut or in the cardiovascular system. In preferred embodiments, the neovascularization in the eye results from age related macular degeneration, diabetic retinopathy, an ocular tumor, central retinal vein occlusion, diabetic macular edema (DME) or pathologic myopia. In a preferred embodiment, the mammal is a human. In an embodiment, the therapeutically effective amount is about 0.01 to about 10 mg/kg body weight, and more preferably, about 0.01 to about 1 mg/kg body weight. In an embodiment, the crystalline form of the dilactate salt is characterized by an X-ray powder diffraction pattern having major diffraction angles. Another aspect of the invention is a process for the preparation of a crystalline form of squalamine dilactate from a non-crystalline form comprising dissolving the non-crystalline squalamine dilactate in a solvent system containing at least two solvents, followed by supersaturating the solvent system until the squalamine dilactate crystallizes from the solvent system. In different embodiments, supersaturation may occur by cooling the solvent system, reducing the volume of the solvent system, adding an additional amount of at least one of the solvents of the at least two solvents or a combination thereof. In a preferred embodiment, at least one solvent of the at least two solvents is 2-propanol, ethanol, water or 2-butanol. Another embodiment of the invention comprises a new method for the production of crystallized squalamine dilactate as part of the manufacturing process that removes the need for a HPLC purification step. BRIEF DESCRIPTION OF THE FIGURES FIG. 1 shows the structure of squalamine. FIG. 2 shows the x-ray diffraction powder pattern for lyophilized squalamine dilactate. FIG. 3 shows a thermogravimetric scan of the lyophilized squalamine dilactate. FIG. 4 shows a Differential Scanning calorimeter profile for lyophilized squalamine dilactate. FIG. 5 shows the crystal structure of squalamine dilactate crystallized from 2-propanol. FIG. 6 shows the x-ray diffraction powder pattern for squalamine dilactate crystallized from 2-propanol. FIG. 7 shows a thermogravimetric scan of the squalamine dilactate crystallized from 2-propanol. FIG. 8 shows a Differential Scanning calorimeter profile for squalamine dilactate crystallized from 2-propanol. FIG. 9 shows the x-ray diffraction powder pattern for squalamine dilactate crystallized from ethanol. FIG. 10 shows a thermogravimetric scan of the squalamine dilactate crystallized from ethanol. FIG. 11 shows a Differential Scanning calorimeter profile for squalamine dilactate crystallized from ethanol. FIG. 12 shows the x-ray diffraction powder pattern for squalamine dilactate crystallized from 2-butanol. FIG. 13 shows a thermogravimetric scan of the squalamine dilactate crystallized from 2-butanol. FIG. 14 shows a Differential Scanning calorimeter profile for squalamine dilactate crystallized from 2-butanol. FIG. 15 shows a scheme depicting a new method for the production of squalamine. FIG. 16 shows a scheme depicting a new method for the production of squalamine dilactate. FIG. 17 shows the x-ray diffraction powder pattern for recrystallized squalamine dilactate produced by the newly described synthesis of squalamine dilactate. DETAILED DESCRIPTION OF THE INVENTION Definitions As used herein, the term “amorphous” refers to a form of a compound that lacks a distinct crystalline structure. As used herein, the term “polymorphic” refers to one of the crystalline forms of a compound or to a compound that has more than one crystalline form. As used herein, the term “organic alcohol” refers to an organic compound with one or more attached hydroxyl groups. As used herein, the term “solvate” refers to a crystalline form of a squalamine that contains solvent molecules as part of the crystal structure. In this case the solvent is not water. As used herein, the term “hydrate” refers to a crystalline form of a squalamine that contains water molecules as part of the crystal structure. As used herein, the term “squalamine” includes the compound shown in FIG. 1 with the chemical name 3β-(N-[3-aminopropyl]-1,4-butanediamine)-7α,24R-dihydroxy-5α-cholestane-24-sulfate. As used herein, the term “aminosterol” refers to a compound with at least one hydroxyl and one amino group directly or indirectly attached to a steroid nucleus. Squalamine is an example of an aminosterol. As used herein, the term “angiogenesis” refers to the formation of new blood vessels, and an angiogenic is a compound that promotes this activity. As used herein, the term “antiangiogenic” refers to the prevention of the formation of new blood vessels or the destruction of newly formed blood vessels, and includes an agent that exhibits one or both of these properties. As used herein, the term “neovascularization” refers to new blood vessel formation in abnormal tissue (as, for example, in a tumor) or in abnormal positions (as, for example, in some conditions of the eye). As used herein, the term “macular degeneration” is intended to encompass all forms of macular degeneration and includes a gradual loss of central vision usually affecting both eyes that occurs especially in the elderly. A slowly progressing form of macular degeneration, usually referred to as the dry form, is marked especially by the accumulation of yellow deposits in the macula lutea and the thinning of the macula lutea. A rapidly progressing form of macular degeneration, usually referred to as the wet form, is marked by scarring produced by bleeding and fluid leakage from new blood vessels formed below the macula lutea. Macular degeneration may exist as either the wet form or the dry form. As used herein, the term “diabetic retinopathy” includes retinal changes occurring in long-term diabetes and is characterized by punctate hemorrhages from newly formed blood vessels in the retina, microaneurysms and sharply defined waxy exudates. As used herein, a “therapeutically effective” amount is an amount of an agent or a combination of two or more agents, which inhibits, totally or partially, the progression of the condition or alleviates, at least partially, one or more symptoms of the condition. A therapeutically effective amount can also be an amount that is prophylactically effective. The amount that is therapeutically effective will depend upon the patient's size and gender, the condition to be treated, the severity of the condition and the result sought. For a given patient, a therapeutically effective amount can be determined by methods known to those of skill in the art. General Squalamine has been shown to exhibit antiangiogenic and antimicrobial properties and is useful for the treatment of diseases associated with the growth of new blood vessels such as solid tumor growth and metastasis, atherosclerosis, age related macular degeneration, diabetic retinopathy, neovascular glaucoma, retinal ischemia, macular edema, inflammatory diseases and the like in an animal, preferably in a mammal and more preferably, in a human. The three basic nitrogen atoms present in the spermidine side chain of squalamine form salts when treated with various acids. One nitrogen atom in the side chain is neutralized by the sulfonic acid at C24 while the other two nitrogen atoms are free to form salts with an added acid. Such squalamine salts include, but are not limited to, dihydrochloride, diacetate, ditrifluoroacetate, digluconate and dilactate. A comparison of various squalamine salts based on their toxicity and stability show the dilactate salt to be a preferred salt. An embodiment of the invention relates to the amorphous dilactate salt form of squalamine. As described below, the dilactate salt can be prepared in an amorphous form through ion exchange chromatography followed by lyophilization or in various crystalline forms by precipitation from different alcoholic solvents. Another aspect of the invention relates to methods for the preparation of the amorphous and the crystalline forms of squalamine dilactate. The complete X-ray structure of the dilactate salt crystallized from 2-propanol has been determined, confirming the stereochemistry at the asymmetric centers of the squalamine molecule as 3β, 5α, 7α and 24 R. Another embodiment of the invention relates to the various crystalline forms of squalamine dilactate. One particular embodiment is the crystalline form of squalamine dilactate precipitated from 2-propanol which is characterized by an X-ray powder diffraction pattern having major diffraction peaks at 12.5, 16.6 and 18.8 degrees. Another particular embodiment relates to the crystalline form of squalamine dilactate precipitated from ethanol which is characterized by an X-ray powder diffraction pattern having major diffraction peaks at 10.2, 13.0 and 16.6 degrees. Another particular embodiment relates to the crystalline form precipitated from 2-butanol which is characterized by an X-ray powder diffraction pattern having major diffraction peaks at 13.1, 17.7 and 18.3 degrees. Another particular embodiment relates to the crystalline form precipitated from ethanol-water which is characterized by an X-ray powder diffraction pattern having major diffraction angles of 12.6, 15.7 and 18.8 degrees. The crystalline forms of squalamine dilactate may exist as solvates, where solvent molecules are incorporated within the crystal structure. As an example, when the solvent contains ethanol, the crystal may contain ethanol molecules. In another embodiment, the solvate may contain an water, and the crystal may be a hydrate containing water in the crystal structure. In another embodiment the crystal may be both a solvate and a hydrate. Another embodiment of the invention comprises a new method for the production of recrystallized squalamine dilactate. This new method utilizes the method described in U.S. Pat. No. 6,262,283 to produce a hydroxy-protected ketosterol 1 (e.g., compound 36 where the protecting group (PG) is —OC(O)-Ph); which is then reacted with azidospermidine to produce the corresponding imine 2; followed by reduction with, for example, NaBH 4 , to produce the azidoaminosterols 3 as a mixture of protected and unprotected 7-alcohols; followed by direct treatment with methanolic potassium hydroxide, followed by hydrogenation in the presence of Raney nickel, to produce crude squalamine. Rather than purification by HPLC and conversion to the dilactate salt by ion exchange chromatography, the crude squalamine is dissolved in ethanol and a two-fold excess of lactic acid is added. The crystalline squalamine dilactate 4 is then precipitated out of solution by the addition of water and, optionally, squalamine dilactate seed crystals. Final purification is then achieved by one or more recrystallizations from aqueous ethanol, preferably containing at least 4% water. This new process produces a better yield and a cleaner product than older methods and results in a considerable cost saving due to the elimination of the HPLC purification step. The squalamine salts of the invention, and in particular, the squalamine dilactate in any of its forms, may be administered alone or as part of a pharmaceutical composition. Pharmaceutical compositions for use in vitro or in vivo in accordance with the present invention may be formulated in a conventional manner using one or more physiologically acceptable carriers comprising excipients and auxiliaries that facilitate processing of the active compounds into preparations which can be used pharmaceutically. Proper formulation is dependent upon the route of administration chosen. Examples of carriers or excipients include, but are not limited to, calcium carbonate, calcium phosphate, various sugars, starches, cellulose derivatives, gelatin and polymers such as polyethylene glycols. One example of a pharmaceutical carrier for the squalamine salts of the invention is a cosolvent system comprising benzyl alcohol, a nonpolar surfactant, a water-miscible organic polymer and an aqueous phase. The proportions of the co-solvent system may be varied considerably without adversely affecting the composition's solubility and toxicity characteristics. Furthermore, the identity of the cosolvent components may be varied: for example, other low-toxicity, nonpolar surfactants may be used instead of polysorbate 80; the fraction size of polyethylene glycol may be varied; and/or other biocompatible polymers may replace polyethylene glycol, e.g., polyvinyl pyrrolidone and sugars or polysaccharides, e.g., dextrose. In addition to carriers, the pharmaceutical compositions of the invention may also optionally include stabilizers, preservatives and/or adjuvants. For examples of typical carriers, stabilizers and adjuvants known to those of skill in the art, see Remington: The Science and Practice of Pharmacy , Lippincott, Williams & Wilkins (2000), which is incorporated by reference in its entirety. Optionally, other therapies known to those of skill in the art may be combined with the administration of the squalamine salts of the invention. More than one aminosterol may be present in a single composition. In vivo administration of squalamine salts of the invention can be effected in one dose, multiple doses, continuously or intermittently throughout the course of treatment. Doses range from about 0.01 mg/kg to about 10 mg/kg, preferably between about 0.01 mg/kg to about 1 mg/kg, and most preferably between about 0.1 mg/kg to about 1 mg/kg in single or divided daily doses. Methods of determining the most effective means and dosages of administration are well known to those of skill in the art and will vary with the composition used for therapy, the purpose of the therapy, the target cell being treated and the subject being treated. Single or multiple administrations can be carried out with the dose level and pattern being selected by the treating physician. Pharmaceutical compositions containing the squalamine salts of the invention can be administered by any suitable route, including oral, rectal, intranasal, topical (including transdermal, aerosol, ocular, buccal and sublingual), parenteral (including subcutaneous, intramuscular, intravenous), intraperitoneal and pulmonary. It will be appreciated that the preferred route will vary with the condition and age of the recipient, and the disease being treated. For treatment of age-related macular degeneration, for example, the preferred routes of administration are topical, subcutaneous, intramuscular and/or intravenous. For oral administration, the squalamine salts of the invention can be formulated readily by combining them with pharmaceutically acceptable carriers well known in the art. Such carriers enable the compounds of the invention to be formulated as tablets, pills, dragees, capsules, liquids, gels, syrups, slurries, suspensions and the like, for oral ingestion by a patient to be treated. Pharmaceutical preparations for oral use can be obtained by combining the active compound with a solid excipient, optionally grinding a resulting mixture, and processing the mixture of granules, after adding suitable auxiliaries, if desired, to obtain tablets or dragee cores. Suitable excipients include, for example, fillers such as sugars, including lactose, sucrose, mannitol, or sorbitol; cellulose preparations such as maize starch, wheat starch, rice starch, potato starch, gelatin, gum tragacanth, methyl cellulose, hydroxypropylmethylcellulose, sodium carboxymethylcellulose and polyvinylpyrrolidone (PVP). If desired, disintegrating agents may be added, such as the cross-linked polyvinyl pyrrolidone, agar, alginic acid or a salt thereof, such as sodium alginate. Pharmaceutical compositions for topical administration of the squalamine salts of the invention may be formulated in conventional ophthalmologically compatible vehicles, such as, for example, an ointment, cream, suspension, lotion, powder, solution, paste, gel, spray, aerosol or oil. These vehicles may contain compatible preservatives such as benzalkonium chloride, surfactants such as polysorbate 80, liposomes or polymers such as methylcellulose, polyvinyl alcohol, polyvinyl pyrrolidone and hyaluronic acid, which may be used for increasing viscosity. For diseases of the eye, preferred topical formulations are ointments, gels, creams or eye drops containing at least one of the aminosterols of the invention. For administration by inhalation, the squalamine salts of the present invention are conveniently delivered in the form of an aerosol spray presentation from pressurized packs or a nebuliser, with the use of a suitable propellant, e.g., dichlorodifluoromethane, trichlorofluoromethane, dichlorotetrafluoroethane, carbon dioxide or other suitable gas. In the case of pressurized aerosol the dosage unit may be determined by providing a valve to deliver a metered amount. Capsules and cartridges of e.g., gelatin for use in an inhaler or insufflator may be formulated containing a powder mix of the compound and a suitable powder base such as lactose or starch. The squalamine salts can be formulated for parenteral administration by injection, e.g., bolus injection or continuous infusion. Formulations for injection may be presented in unit dosage form, e.g., in ampoules or in multi-dose containers, with an added preservative. The compositions may take such forms as suspensions, solutions or emulsions in oily or aqueous vehicles, and may contain formulatory agents such as buffers, bacteriostats, suspending agents, stabilizing agents, thickening agents, dispersing agents or mixtures thereof. Pharmaceutical formulations for parenteral administration include aqueous solutions of the active compounds in water-soluble form. Additionally, suspensions of the active compounds may be prepared as appropriate oily injection suspensions. Suitable lipophilic solvents or vehicles include fatty oils such as sesame oil, or synthetic fatty acid esters, such as ethyl oleate or triglycerides or liposomes. Aqueous injection suspensions may contain substances that increase the viscosity of the suspension, such as sodium carboxymethyl cellulose, sorbitol or dextran. Optionally, the suspension may also contain suitable stabilizers or agents that increase the solubility of the compounds to allow for the preparation of highly concentrated solutions. In a preferred embodiment, the squalamine salts of the invention are dissolved in a 5% sugar solution, such as dextrose, before being administered parenterally. For injection, the squalamine salts of the invention may be formulated in aqueous solutions, preferably in physiologically compatible buffers such as Hanks's solution, Ringer's solution or physiological saline buffer. For transmucosal administration, penetrants appropriate to the barrier to be permeated are used in the formulation. Such penetrants are generally known in the art. The squalamine salts may also be formulated in rectal compositions such as suppositories or retention enemas, e.g., containing conventional suppository bases such as cocoa butter or other glycerides. The squalamine salts may also be combined with at least one additional therapeutic agent. Exemplary agents include, for example, anticancer, antibiotic, antiviral, antiangiogenic or another treatment for neovascularization in the eye. Without further description, it is believed that one of ordinary skill in the art can, using the preceding description and the following illustrative examples, make and utilize the compounds of the present invention and practice the claimed methods. The following working examples therefore, specifically point out preferred embodiments of the present invention, and are not to be construed as limiting in any way the remainder of the disclosure. EXAMPLES Example 1 Preparation of Amorphous Squalamine Dilactate Crude squalamine was prepared according to the methods described in U.S. Pat. No. 6,262,283, U.S. Pat. No. 6,610,866 and U.S. patent application Ser. No. 10/268,660. The crude squalamine was dissolved in water, acidified with trifluoroacetic acid (TFA) and then purified by reverse phase HPLC using a C 18 YMC ODS-AQ column or equivalent and a binary solvent system. The HPLC chromatography was performed to collect fractions of product that meet the drug substance specifications. The fractions of pure squalamine TFA salt were concentrated prior to salt conversion. Conversion of the squalamine TFA salt to squalamine dilactate salt was accomplished by adsorption of the TFA salt to Amberchrom resin or its equivalent. The resin was then washed extensively with 1% acetonitrile in water, sodium bicarbonate and 1% acetonitrile in water; and finally with an excess of L-(+) lactic acid dissolved in water. The dilactate salt of squalamine was eluted with a stepwise increase in the percentage of acetonitrile in water. The fractions containing squalamine dilactate were pooled, concentrated and lyophilized. Analysis of the material for lactic acid and squalamine showed a ratio of two moles of lactic acid per mole of squalamine. The characterization of the lyophilized squalamine dilactate is described below. X-Ray Diffraction Powder Pattern The powder x-ray diffraction scans on a sample of lyophilized squalamine dilactate were performed from 4.0 to 45.0 degrees (2 in compliance with USP Method <941>) while a polycarbonate film covered the sample. The pertinent data is shown in FIG. 2 and summarized in the Table below. These data show an amorphous halo with a few discrete peaks indicating a low or partial crystallinity. Angle (° theta-2 theta) Sample Preparation 15.5-15.6 17.3-17.5 21.3-21.5 Lyophilized Squalamine 286 391 107 Thermogravimetric Analysis Thermogravimetric analysis involves the determination of the weight of a specimen as a function of temperature as per USP <891>. The samples of lyophilized squalamine dilactate were prepared in a nitrogen atmosphere in a humidity-controlled glovebox. Analyses were completed using a Perkin Elmer TGA7 with TAC 7/DX Instrument Controller and Pyris Software Version 4.01. Nitrogen, NF was used at a flow rate of 20 mL/minute. The samples were warmed at a controlled rate of 10° C. per minute to generate better sensitivity and at 2° C. per minute in order to acquire a better resolution to a final temperature of 180° C. The results for the 2° C. per minute scan are shown in FIG. 3 and the data summarized in the table below. The data exhibit a single weight loss of 2.32% and degradation onset at a temperature of 136.9° C. Squalamine Primary Peak of Dilactate Rate Mass Secondary Secondary Onset to Preparation (° C./min) Loss Mass Loss Mass Loss Degradation Lyophilized 2 2.32% N/A N/A 136.9° C. Squalamine Differential Scanning Calorimetry Samples were analyzed by high temperature differential scanning calorimetry and were run at 2° C. and 10° C. per minute. Thermal transitions acquired during a scan rate of 2° C. per minute are considered to be more accurate and are the calculations reflected in the conclusion. All events listed are endothermic peak temperatures unless otherwise noted. Examples of additional events include an “Exo” indicating an exothermic event or “Tg′” which indicates a phase transition. The lack of a notable thermal event on a particular scan is indicated by “none”. During analysis of lyophilized squalamine dilactate, an exothermic event was detected at an onset temperature of 52.7° C. during a scan at 2° C. per minute. A phase transition (Tg) event, occurring at a temperature of 62.0° C. during a scan at 10° C. per minute, did not have a corresponding thermal event when scanned at 2° C. per minute. A phase transition is indicative of an amorphous portion of the dried material softening and changing structure. Two endothermic events were observed at peak temperatures of 127.7° C. and 157.7° C. during the scan at 2° C. per minute. The largest change in specific heat associated with an endothermic event for lyophilized squalamine dilactate was a change in specific heat of 8.15 J/g which was observed during the 10° C. per minute scan at a temperature of 166.51° C. The change in specific heat associated with an endothermic event is correlated to the amount of energy required to melt that material. The results of the 2° C. per minute scan are shown in FIG. 4 and summarized in the table below. Squalamine Dilactate Rate Preparation (° C./min) 1 2 3 4 5 Lyophilized 2 52.7 None 127.7 157.7 none Squalamine (Exo) 10 62.0 71.3  98.3 130.1 166.5 (Tg) (Exo) (Exo) Example 2 A Study of Local Irritancy of 5-Day Repeated Intravenous Injections of Different Salt Forms (Ditrifluoroacetate, Dilactate, Digluconate, Diacetate) of MSI-1256 (Squalamine) in Mice Summary: Five-day repeated injections of various salt forms of squalamine (2.5 mg/kg/day) caused swelling, bruising and irritation of the mouse-tails. Treatment with squalamine dilactate and squalamine digluconate, was tolerated slightly better than treatment with squalamine diacetate and squalamine ditrifluoroacetate, although swelling, bruising and irritation were observed with all injected salt forms of squalamine administered repeatedly at a dose of 2.5 mg/kg/day using 0.25 mg/mL solutions. Objective: To determine the local irritancy of 5-day repeated daily intravenous doses of squalamine salt forms in tails of CD-1®BR mice. Material and Methods: (Animals): Forty-eight male CD-1®BR mice (Charles River Lab). Mean body weight at study initiation was 20.6 gm. (Housing Environment): Mice were housed as groups (maximum 10 mice/box) in plastic mouse boxes with hardwood chip bedding and wire lids. They had access to food (Purina Mouse Chow) and water in bottles ad lib. The boxes were housed in isolator racks that were supplied with one-pass through filtered room air. The room in which the animals were housed was on a 12-hour on/12-hour off light cycle and had controlled temperature (range: 67-76° F.) and humidity (range: 40-70% relative humidity). Test Articles: squalamine ditrifluoroacetate, 69.7% active moiety squalamine diacetate, 80.0% active moiety squalamine dilactate, 76.0% active moiety squalamine digluconate, 55.0% active moiety Magainin-2-amide, positive control Vehicle: 5% Dextrose in Water (D5W) (Baxter I.V. bag, sterile) Solution prep: A 0.36 mg/mL solution of squalamine ditrifluoroacetate salt (equivalent to 0.25 mg/mL squalamine ditrifluroacetate active moiety) was prepared in D5W. A 0.31 mg/mL solution of squalamine diacetate salt (equivalent to 0.25 mg/mL squalamine diacetate active moiety) was prepared in D5W. A 0.33 mg/mL solution of squalamine dilactate salt (equivalent to 0.25 mg/mL squalamine dilactate active moiety) was prepared in D5W. A 0.45 mg/mL solution of squalamine digluconate salt (equivalent to 0.25 mg/mL squalamine digluconate active moiety) was prepared in D5W. A 1.0 mg/mL solution of magainin-2-amide (positive control) was prepared in D5W. Protocol: Mice were randomly assigned to six groups (8 mice/group) and received daily intravenous (i.v.) injections of solutions of D5W or squalamine salts in the tail vein for five days (study days 0, 1, 2, 3, and 4). Injectate volumes of 10 mL/kg body weight using D5W or a 0.25 mg/mL solution of squalamine salts resulted in doses of 0 mg/kg/day for the mice in D5W group (Group 1) and 2.5 mg/kg/day of squalamine salt active moiety for all squalamine salt treated mice. One group of eight mice received test article magainin-2-amide (10 mg/kg/day; 10 mL/kg/day of a 1 mg/mL solution), which was previously determined to be local vein irritant, as a positive control. Mice were not injected with test article if severity of bruising or swelling (edema) warranted the omission of the injection. Survival was monitored and clinical signs were observed daily for five days of administration of squalamine salts and four days after the last dose (Study Day 8). Clinical signs of irritancy, edema, and bruising were made on days 1, 2, 3, and 4 approximately 24-hours after each injection and immediately prior to that day's injection. The pHs of the solutions of all test articles except magainin-2-amide were checked on Study Day 3. (The same solutions were used throughout the study so a similar pH on all study days may be assumed). Results: Animals that were administered repeated intravenous (i.v.) doses of various salt forms of squalamine in the tail vein had irritated, swollen (edema) and/or bruised tails by Study Day 2. The number of mouse tails which were bruised and edematous as well as the severity of bruising and edema was directly related to the number of injections. To assess recovery, tails were observed on Study Day 8, which was four days post-last injection. On Day 8, the tails of mice in Groups 3 and 4 (having received the dilactate and digluconate salt forms, respectively) were slightly irritated and bruised. The tails of mice in Groups 5 and 6 (having received the diacetate and ditrifluoroacetate salt forms, respectively) were similarly irritated and bruised. One (⅛) mouse tail in Group 5 was necrotic. In Group 6, one (⅛) mouse tail fell off, and one (⅛) mouse tail was also necrotic. Group 2 (positive control) mice showed slight or moderate edema in tails during dosing (⅛, ⅜, ⅝, and ⅝ on Days 1, 2, 3 and 4) and recovered by Day 8. The pHs of all solutions of squalamine salt forms were approximately 6. Conclusions: Repeated injections of all salt forms of squalamine caused swelling (edema), bruising and irritation of the mouse tails. The clinical symptoms suggest that treatment with squalamine dilactate and squalamine digluconate was tolerated better than treatment with squalamine diacetate and squalamine ditrifluoroacetate. Thus, an unexpected advantage of the squalamine dilactate and squalamine digluconate salts over other tested squalamine salts is less venous irritation, i.e., less toxicity, experienced by the recipient, especially at the site of intravenous administration. Example 3 Accelerated Stability Study of Four Salt Forms of Squalamine An accelerated stability study (temperature-based) lasting four weeks was performed on squalamine in four different salt forms. The four salt forms were: dihydrochloride, diacetate, dilactate and D-digluconate. The samples were subjected to temperatures of 40° C., 60° C. and 80° C. The following table summarizes the results of % purity of main peak, squalamine, based on total integrated area. The analysis was performed using reversed-phase HPLC of o-phthaldialdehyde derivatized samples. TABLE 1 Salt Form T = 0 4 w 40° C. 4 w 60° C. 4 w 80° C. Dihydrochloride 90.7% 84.3% 85.3% 82.4% Diacetate 94.4% 87.3% 81.5% 62.9% Dilactate 91.8% 80.9% 80.6% 70.9% Digluconate 87.8% 72.6% 60.7% 4.9% Table 1 shows how each salt form has degraded over time at elevated temperatures. The results of this stability study indicate that squalamine dilactate is surprisingly stable under increasingly severe conditions, especially compared to the diacetate and digluconate salts. This advantageous stability of the squalamine dilactate salt coupled with its low toxicity (as shown in Example 2) were important factors in selecting the squalamine dilactate salt form for further development. Example 4 Preparation of Crystalline Squalamine Dilactate from 2-Propanol A supersaturated solution of amorphous squalamine dilactate was produced by heating an excess of squalamine dilactate to 90° C. in a mixture of 10 ml of 2-propanol plus 100 μl of water. The excess squalamine dilactate was filtered off and the solution was cooled to −20° C. A precipitate of white needles formed, the supernatant was removed and the solid dried in a vacuum desiccator. The resulting crystalline material was observed to be non hygroscopic as it did not gain weight when left at room temperature uncovered for one hour. Single Crystal X-Ray Diffraction Pattern Determination Single crystals suitable for X-ray study were obtained from a solution of 2-propanol and water as described above. The biggest crystal, with the dimensions of 0.025, 0.10, and 1.10 mm was chosen for the study. The crystal was mounted on a Nonius Kappa CCD instrument with molybdenum radiation and CCD area detector. The crystal was cooled to 173° K using a nitrogen stream cooled by liquid nitrogen. The preliminary measurements showed that the diffraction was very weak beyond 22 degrees in theta and that the crystal belonged to the monoclinic space group. To enhance the differentiation of chiral isomers, it was decided that the data would be collected in lower crystal system (namely triclinic). The data were collected by exposing the crystal for 500 seconds per degree of crystal rotation. The total data collection took 32 hours. The data were processed to obtain the final intensity of the diffraction pattern and all the unique measurements were kept separate without applying the Friedel law. The space group analysis showed that there was no systematic absence. The diffraction pattern analysis showed that the crystal belonged to a non-centric space group and a possible two fold symmetry along the b-axis, suggesting that it may belong to P2 space group. All these were consistent with what is expected for a chiral molecule, which cannot belong to centric space group and cannot have mirror or glide symmetries. The intensity analysis showed that the data were becoming weak at the higher angles of theta. The average intensity at the 21-22 degree theta range was just 1.9 times of the average background. However, there were enough strong data to provide the molecular structure with proper absolute configuration (chirality). The structure was readily solved by direct methods in the space group P2 suggested by the data analyses. Refinement of the structure by least-squares followed by difference Fourier showed the presence of solvent molecules. Many water molecules and one disordered 2-propanol molecule were detected. The occupancies of the molecules were verified by refinement and one water molecule was found to have only half (50%) occupancy. All the non-hydrogen atoms were refined with anisotropic displacement parameters (ADP). Hydrogen atoms connected to carbon and nitrogen atoms were included at calculated positions. For hydroxyl groups and water molecules, only when reasonable atoms were found in difference Fourier, they were included. For the disordered 2-propanol molecule and for oxygen for which no reasonable atom could be located from different Fourier were not included in the calculation. The refinement used 6119 intensity data and refined 612 parameters. The final residual factor (R-factor) was 0.087, which unambiguously proves the structure of the molecule. There are several polar (electron-deficient) hydrogens at cationic nitrogens and hydroxyl oxygens. There are also several electron rich oxygens at anionic centers. This leads to a network of hydrogen bonding formation. Also, many water molecules join the hydrogen-bonding network. The detailed are described below. H-Donor H-Acceptor Distance (A) Symmetry O1 O2W 2.692 x, y, z N1(H1A) O10 2.835 x − 1, y − 1, z N1(H1B) O9 2.842 x − 1, y − 2, z N1(H1B) O11 2.854 x − 1, y − 2, z N2(H2A) O6 2.734 −x, y, 2 − z N2(H2B) O7 2.876 −x, y − 1, 2 − z N2(H2B) O8 2.835 −x, y − 1, 2 − z N3(H3C) O3 2.886 −x, y, 2 − z N3(H3A) O1W 2.804 x − 1, y − 1, z + 1 N3(H3B) O1S or O1S′ 2.940 or 2.945 x − 1, y − 1, z + 1 O2W(H3W) O7 2.811 x, y, z O2W(H4W) O9 2.762 1 − x, y − 1, 2 − z O8(H8) O3W 2.747 x, y, z O3W O4W 2.777 x, y, z O3W O6 2.665 x, 1 + y, z O4W O10 2.749 x, y, z O4W O4W 2.840 1 − x, y, 2 − z O11 O4W 2.779 x, 1 + y, z O1S or O1S′ O3W 2.730 or 2.857 1 − x, y, 1 − z The distances are given between the non-hydrogen atoms and where available the hydrogens through which the bonding formed are shown in parenthesis. The crystal structure of squalamine dilactate is shown in FIG. 5 . A unit cell was determined from the single crystal X-ray data of the hydrated form. It was monoclinic with P2 symmetry, Z=2, and the following dimensions: a=19.3999 Å, b=6.5443 Å, c=20.9854 Å, alpha=gamma=90°, beta=92.182° and V=2662.3 Å 3 . X-Ray Diffraction Powder Pattern The powder x-ray diffraction scans on a sample of squalamine dilactate crystallized from 2-propanol were performed from 4.0 to 45.0 degrees (2 in compliance with USP Method <941>) while the sample was covered by a polycarbonate film. The pertinent data, consisting of distinct crystalline peaks is shown in FIG. 6 and indicates a crystalline material. Angle (° theta-2 theta) Sample Preparation 12.5 16.6 18.8 Crystallized from 2- 890 829 756 propanol/water Thermogravimetric Analysis Thermogravimetric analysis involves the determination of the weight of a specimen as a function of temperature as per USP <891>. The samples were prepared in a nitrogen atmosphere in a humidity controlled glovebox. Analyses were completed using a Perkin Elmer TGA7 with TAC 7/DX Instrument Controller and Pyris Software Version 4.01. Nitrogen, NF was used at a flow rate of 20 mL/minute. The samples were warmed at a controlled rate of 10° C. per minute to generate better sensitivity and at 2° C. per minute in order to acquire a better resolution to a final temperature of 180° C. This crystallized material had two distinct volatile weight loss events. The initial event yielded a 1.38% weight loss. The second event yielded an average weight loss of 1.54% with a peak event observed at a temperature of 103.6° C. when tested at 2° C. per minute. The total weight loss incurred by the sample was 2.92%. The two distinct weight loss events suggest that a bound form of water existed within the sample matrix. The initial loss is most likely due to the driving off of volatile constituents adsorbed to the material surface. The second weight loss event occurred due to a release of absorbed water associated with the sample matrix, which was most likely a crystalline hydrate, at a peak temperature of 103.6° C. The second distinct release of moisture from the sample at a specific temperature suggests a breakdown of a large portion of crystalline material. The results for the 2° C. per minute scan are shown in FIG. 7 and the data summarized in the table below. Squalamine Primary Peak of Dilactate Rate Mass Secondary Secondary Onset to Preparation (° C./min) Loss Mass Loss Mass Loss Degradation Crystallized 2 1.38% 103.6° C. 1.54% 130.8° C. from 2-propanol/ water Differential Scanning Calorimetry Samples were analyzed by high temperature differential scanning calorimetry and were run at 2° and 10° C. per minute. Thermal transitions acquired during a scan rate of 2° C. per minute are considered to be more accurate and are the calculations reflected in the conclusion. All events listed are endothermic peak temperatures unless otherwise noted. Examples of additional events include an “Exo” indicating an exothermic event or “Tg′” which indicates a phase transition. The lack of a notable thermal event on a particular scan is indicated by “none”. The first thermal event characterized by DSC at 2° C. per minute for this crystallized material was an endothermic event that occurred at a temperature of 73.6° C. Endothermic events are attributable to the initial melt of a crystallized material. The most significant thermal event was an additional endothermic event that occurred at a temperature of 107.3° C. This endothermic event coincides with the peak weight loss temperature of 103.6° C. during the TGA scan of this particular material. Three additional exothermic events occurred at temperatures of 126.6°, 157.3°, and 164.1° C. The results of the 2° C. per minute scan are shown in FIG. 8 and summarized in the table below. Squalamine Dilactate Rate Preparation (° C./min) 1 2 3 4 5 Crystallized from 2 73.6 107.3 126.6 157.3 164.1 2-propanol/water 10 79.8 112.8 130.5 none 165.6 Example 5 Preparation of Crystalline Squalamine Dilactate from Ethanol A supersaturated solution of amorphous squalamine dilactate was produced by heating an excess of squalamine dilactate to 90° C. in a mixture of 10 ml of ethanol plus 100 μl of water. The excess squalamine dilactate was filtered off and the solution was cooled to −20° C. A precipitate of white needles formed, the supernatant was removed and the solid dried in a vacuum desiccator. X-ray Diffraction Powder Pattern The powder x-ray diffraction scans on a sample of squalamine dilactate crystallized from ethanol were performed from 4.0 to 45.0 degrees (2 in compliance with USP Method <941>) while the sample was covered by a polycarbonate film. The pertinent data, consisting of distinct crystalline peaks is shown in FIG. 9 and summarized in the table below and indicates a crystalline material. Angle (° theta-2 theta) Sample Preparation 10.2 13.0 16.6 Crystallized from ethanol/water 1826 2305 1817 Thermogravimetric Analysis Thermogravimetric analysis involves the determination of the weight of a specimen as a function of temperature as per USP <891>. The samples were prepared in a nitrogen atmosphere in a humidity controlled glovebox. Analyses were completed using a Perkin Elmer TGA7 with TAC 7/DX Instrument Controller and Pyris Software Version 4.01. Nitrogen, NF was used at a flow rate of 20 mL/minute. The samples were warmed at a controlled rate of 10° C. per minute to generate better sensitivity and at 2° C. per minute in order to acquire a better resolution to a final temperature of 180° C. This squalamine dilactate sample had two distinct volatile weight loss events. The initial event yielded a 2.99% weight loss. The second event yielded an average weight loss of 1.49% with a peak event observed at a temperature of 106.7° C. when tested at 2° C. per minute. The total weight loss incurred by the sample was 4.48%. The two distinct weight loss events suggest that a bound form of water existed within the sample matrix. The initial loss is most likely due to the driving off of volatile constituents adsorbed to the material surface. The second weight loss event occurred due to a release of absorbed water associated with the sample matrix, which was most likely a crystalline hydrate, at a peak temperature of 106.7° C. The second distinct release of moisture from the sample at a specific temperature suggests a breakdown of a large portion of crystalline material. Temperatures relating to the onset of degradation were determined as well using TGA. The average onset to degradation temperature for this sample material was 152.3° C. The results for the 2° C. per minute scan are shown in FIG. 10 and the data summarized in the table below. Squalamine Rate Primary Peak of Dilactate (° C./ Mass Secondary Secondary Onset to Preparation min) Loss Mass Loss Mass Loss Degradation Crystallized from 2 2.99% 106.7° C. 1.49% 152.3° C. ethanol/water Differential Scanning Calorimetry Samples were analyzed by high temperature differential scanning calorimetry and were run at 2° C. and 10° C. per minute. Thermal transitions acquired during a scan rate of 2° C. per minute are considered to be more accurate and are the calculations reflected in the conclusion. All events listed are endothermic peak temperatures unless otherwise noted. Examples of additional events include an “Exo” indicating an exothermic event or “Tg′” which indicates a phase transition. The lack of a notable thermal event on a particular scan is indicated by “none”. The first and most significant thermal event characterized by DSC at 2° C. per minute for this material, was an endothermic event that occurred at a temperature of 112.3° C. Endothermic events are attributable to the initial melt of a crystallized material. This endothermic event coincides with the peak weight loss temperature of 106.7° C. during the TGA scan of this particular material. Two additional exothermic events were detected at temperatures of 144.7° C. and 175.4° C. The endothermic events, occurring at the temperatures of 141.1° C. and 151.6° C. during a scan at 10° C. per minute, did not have a corresponding thermal event when scanned at 2° C. per minute. The thermal events observed during the scan of the material indicate the melt of crystalline material. The most significant endothermic event observed during the 2° C. per minute scan of the material, occurring at a peak temperature of 112.4° C., resulted in a change in specific heat of 18.16 J/g. The change in specific heat associated with an endothermic event is correlated to the amount of energy required to melt that material. Therefore, the endothermic event, occurring the temperature of 112.4° C. is considered to represent the most stable crystalline material present. The results of the 2° C. per minute scan are shown in FIG. 11 and summarized in the table below. Squalamine Dilactate Rate Preparation (° C./min) 1 2 3 4 5 Crystallized from 2 112.3 None 144.7 None 175.4 ethanol/water 10 114.6 141.1 147.4 151.6 173.1 Example 6 Stability of Crystalline Squalamine Dilactate Precipitated from Ethanol Samples of amorphous squalamine dilactate and squalamine dilactate recrystallized from ethanol as described in Example 5 were placed in scintillation vials and heated in an oven at 70° C. for three days. A portion of each heat stressed sample was then analyzed by HPLC with ELSD detection and the results compared with the HPLC analysis of unstressed material. The result of the HPLC analysis is shown in the table below. 24 S Un- Material Isomer Squalamine Lactylamide Unknown known Amorphous 0.833% 98.978% 0.115% 0.054% 0.020% Squalamine Unstressed Amorphous 0.864% 98.530% 0.519% 0.052% 0.034% Squalamine Heat Stressed Crystallized 0.715% 99.187% 0.070% 0.032% <0.02% Squalamine Unstressed Crystallized 0.714% 99.168% 0.087% 0.031% <0.02% Squalamine Heat Stressed The result demonstrates a significant increase in the percentage of the lactylamide impurity when the amorphous squalamine is heat stressed but no significant increase in the crystallized material. We therefore conclude that recrystallization of squalamine dilactate is a method for improving the stability of the material. This improved stability is advantageous in the preparation and storage of the crystalline squalamine dilactate salt and its various formulations. Example 7 Preparation of Crystalline Squalamine Dilactate from 2-Butanol A supersaturated solution of amorphous squalamine dilactate was produced by heating an excess of squalamine dilactate to 90° C. in a mixture of 10 ml of 2-butanol plus 100 μl of water. The excess squalamine dilactate was filtered off and the solution was cooled to −20° C. A precipitate of white needles formed, the supernatant was removed and the solid dried in a vacuum desiccator. X-Ray Diffraction Powder Pattern The powder x-ray diffraction scans on a sample of squalamine dilactate crystallized from 2-butanol were performed from 4.0 to 45.0 degrees (2 in compliance with USP Method <941>) while the sample was covered by a polycarbonate film. The pertinent data, consisting of distinct crystalline peaks is shown in FIG. 12 and summarized in the table below and indicates a crystalline material. Angle (° theta-2 theta) Sample Preparation 13.1 17.7 18.3 Crystallized from 2-butanol/ 939 937 967 water Thermogravimetric Analysis Thermogravimetric analysis involves the determination of the weight of a specimen as a function of temperature as per USP <891>. The samples were prepared in a nitrogen atmosphere in a humidity controlled glovebox. Analyses were completed using a Perkin Elmer TGA7 with TAC 7/DX Instrument Controller and Pyris Software Version 4.01. Nitrogen, NF was used at a flow rate of 20 mL/minute. The samples were warmed at a controlled rate of 10° C. per minute to generate better sensitivity and at 2° C. per minute in order to acquire a better resolution to a final temperature of 180° C. This squalamine dilactate sample had two distinct volatile weight loss events. The initial event yielded a 2.69% weight loss. The second event yielded an average weight loss of 3.34% with a peak event observed at a temperature of 101.8° C. when tested at 2° C. per minute. The total weight loss incurred by the sample was 6.03%. The two distinct weight loss events suggest that a bound form of solvent existed within the sample matrix. The initial loss is most likely due to the driving off of volatile constituents adsorbed to the material surface. The second weight loss event occurred due to a release of absorbed water associated with the sample matrix, which was most likely a crystalline hydrate, at a peak temperature of 101.8° C. The second distinct release of moisture from the sample at a specific temperature suggests a breakdown of a large portion of crystalline material. Temperatures relating to the onset of degradation were determined as well using TGA. The average onset to degradation temperature for this sample material was 145.0° C. The results for the 2° C. per minute scan are shown in FIG. 13 and the data summarized in the table below. Squalamine Rate Primary Peak of Dilactate (° C./ Mass Secondary Secondary Onset to Preparation min) Loss Mass Loss Mass Loss Degradation Crystallized from 2 2.69% 101.8° C. 3.34% 145.0° C. 2-butanol/water Differential Scanning Calorimetry Samples were analyzed by high temperature differential scanning calorimetry and were run at 2° and 10° C. per minute. Thermal transitions acquired during a scan rate of 2° C. per minute are considered to be more accurate and are the calculations reflected in the conclusion. All events listed are endothermic peak temperatures unless otherwise noted. Examples of additional events include an “Exo” indicating an exothermic event or “Tg′” which indicates a phase transition. The lack of a notable thermal event on a particular scan is indicated by “none”. The first thermal event characterized by DSC at 2° C. per minute for this material, was a glass transition (Tg) event that occurred at a temperature of 45.7° C. A glass transition event is often attributed to some amount of amorphous material. An endothermic event was detected at a temperature of 100.2° C., resulting in a change in specific heat of 19.99 J/g. Endothermic events are attributable to the initial melt of a crystallized material. The change in specific heat associated with an endothermic event is correlated to the amount of energy required to melt that material. This endothermic event coincides with the peak weight loss temperature of 101.8° C. during the TGA scan of this particular material. Therefore, the endothermic event, occurring the temperature of 100.2° C. is considered to represent the most stable crystalline material present. Two additional exothermic events were detected at temperatures of 146.4° and 177.1° C. The endothermic event, occurring at the temperatures of 153.8° during a scan at 10° C. per minute, did not have a corresponding thermal event when scanned at 2° C. per minute. The results of the 2° C. per minute scan are shown in FIG. 14 and summarized in the table below. Squalamine Dilactate Rate Preparation (° C./min) 1 2 3 4 5 Crystallized from 2 45.7 100.2 146.4 None 177.1 2-butanol/water (Tg) 10 42.2 103.2 146.1 153.8 179.8 (Tg) (onset) Example 8 Improved Method for the Manufacturing of Crystalline Squalamine Dilactate Compound 36 was prepared according to the methods described in U.S. Pat. Nos. 6,262,283 and 6,610,866 and U.S. patent application Ser. No. 10/268,660. Approximately 490.0 gms. (2.0 Moles) of azidospermidine dihydrochloride was dissolved in 22.5 L of pyridine at ambient temperature. Approximately 8.0 L (4.0 Moles) of a 0.5 M solution of sodium methoxide-methanol solution was added and the mixture was stirred for about 0.5 hours. Approximately 641.0 gms. (1.0 Mole) of compound 36 was added and the reaction mixture stirred for an additional two hours. The reaction mixture was concentrated to dryness in vacuo (max. 43° C./171 mbar) to remove water, approximately 11.28 L of pyridine was added and the solvent was again distilled off in vacuo. (max. 43° C./171 mbar). Approximately 22.5 L of methanol was added and the obtained suspension cooled to less than about −75° C. Approximately 114 gms. (3.0 Moles) of sodium borohydride was added and the reaction mixture was stirred at less than about −75° C. until compound 36 was transformed as analyzed by HPLC. The mixture was heated to about 15 to about 25° C. and then 2.7 L of distilled water was added to the solution. The solution was concentrated at reduced pressure and a temperature of less than 65° C. to a final volume of about 26.8 L. Approximately 13.4 L of 2-Butanol was added and the mixture stirred before allowing the layers to separate. The lower aqueous layer was removed for disposal at the completion of the batch. (If there is no separation, add MTBE (up to 5 L) to the mixture to aid in separating layers.) The organic phase was washed with 2.7 L of distilled water, the aqueous phase back washed with 17.2 L of 2-Butanol and the two 2-Butanol phases combined. The organic portion, crude compound 40, was concentrated to dryness in vacuo to be used for the preparation of squalamine without further purification. Approximately 796.18 gms. (1.0 Mole) of crude compound 40 was dissolved in 5.7 L of methanol and approximately 280 gms. (5.0 Moles) of potassium hydroxide was added. The reaction mixture was heated at reflux (about 64° C.) until all of the Compound 40 was consumed. Approximately 198 gms. of Raney Nickel catalyst was added and the reaction mixture was hydrogenated at a temperature of 15-25° C. under 2-3 bars of hydrogen pressure until Compound 38 was consumed as analyzed by TLC. The reaction mixture was filtered to remove the catalyst using Celite 545 as a filter aid. The filter cake was washed two times with 800 mL methanol and the combined filtrate and washes were concentrated in vacuo at a temperature of less than about 60° C. to a volume of 6.7 L. Approximately 18.8 L of 2-Butanol was added to the concentrated solution and the solution concentrated under reduced pressure at less than about 60° C. to about 5.36 L. Approximately 13.4 L of methyl t-butyl ether was added and the solution cooled to less than about −5° C. The precipitated product was collected, the filter cake washed two times with 1.3 L of methyl t-butyl ether and the product dried under vacuum at about 25 to about 35° C. A total of 490 gms. of crude squalamine was obtained representing a yield of 75.5%. The synthesis scheme for crude squalamine is shown in FIG. 15 . Recrystallization: Approximately 650 gms. (1.0 Mole) of crude squalamine was mixed into 11.05 L of ethanol to form a cloudy solution. The solution was filtered through a filter coated with filter aid and the filter cake washed with 650 ml. of ethanol. Approximately 494 ml of water and approximately 360.3 gms. (4.0 Moles) of L-(+) Lactic Acid was added to the filtered solution with stirring. The resulting solution was filtered through a 0.22 gm filter and the container and filter washed with 650 ml. of ethanol. The filtrate was cooled to about 0 to about 5° C. for at least 12 hours without stirring and then approximately 100 mg of recrystallized squalamine dilactate seed crystals were added. The solution was maintained at about 0 to about 5° C. without stirring for at least 48 more hours and then the resulting precipitation was agitated at less than about 5° C. to form a homogeneous suspension. The solids were collected and the filter cake washed with 650 ml. of cold (0 to 5° C.) ethanol. The product was dried in vacuo at about 40° C. (±2° C.) to yield a total of 614 g (76.0% yield) of crystallized squalamine dilactate. The synthesis scheme for crude squalamine dilactate is shown in FIG. 16 . Approximately 1 kg. of crystallized or recrystallized squalamine dilactate was combined with 18 L. of ethanol and 760 ml. of water. The suspension was heated to about 40 to about 50° C. with stirring to form a solution and then filtered through a 0.22 μm filter. The container and filter were washed with 1 L. of ethanol and the total filtrate cooled to 20° C. (±2° C.) for at least about twelve hours. Approximately 100 mg of recrystallized squalamine dilactate seed crystals were added the solution was maintained at about 20° C. (±2° C.) without stirring for at least 48 more hours. The resulting precipitation was agitated to form a homogeneous suspension and the solids collected. The filter cake was washed with 1.0 L. of cold (0 to 5° C.) ethanol and the product dried in vacuo at about 40° C. (±2° C.) to yield a total of 900 gms (90.0% yield) of recrystallized squalamine di lactate. The analysis of the crystalline squalamine dilactate produced by this process is shown in the table below. Test Specification Results HPLC Purity >95.00%   99.0% 24-S ≦1.7% 0.74% 3-α ≦0.1% <0.1% Lactyl Amide  <1.5% 0.17% Des-Sulfate ≦0.1% <0.1% Lactic Acid   <30% 22.89%  Water by Karl Fischer   <10% 2.18% HPLC MS 628 ± 1amu Conforms NMR Conforms to ref. Conforms FTIR Conforms to ref. Conforms XRD No specification Completed DSC No specification Mp 143.9° Purity 99.99% Residual solvents (Ethanol) 5000 ppm <200 ppm Sodium No specification 80.5 mg/kg Potassium No specification  520 mg/kg The X-ray diffraction powder pattern, which was determined as described in Example 1 above, is shown in FIG. 17 and the intensity of the major peaks is shown in the table below. Angle (° theta-2 theta) Sample Preparation 12.6 15.7 18.8 Crystallized from ethanol/water 977 891 1333 in the squalamine manufacturing process The powder pattern indicates that even though the squalamine dilactate was recrystallized from ethanol/water as in Example 5, a different polymorphic form has been produced (compare FIGS. 9 and 17 ). This is likely due to the use of 4% water in the manufacturing process as opposed to 1% water in Example 5 and the fact that the material was crystallized at 20° C. instead of −20° C. There is also evidence from the Karl Fisher titration that the recrystallized material from the manufacturing process is a monohydrate. This new manufacturing process also produces a better yield and a purer product than the process described in U.S. Pat. No. 6,262,283. While the invention has been described and illustrated herein by references to various specific materials, procedures and examples, it is understood that the invention is not restricted to the particular combinations of material and procedures selected for that purpose. Numerous variations of such details can be implied as will be appreciated by those skilled in the art. All patents, patent applications and other references cited throughout this application are herein incorporated by reference in their entirety.

The invention relates to select squalamine salts, methods of their synthesis, their therapeutic use and their advantages relating to manufacturing, product stability and toxicity. More specifically, this application is directed to various forms of the dilactate salt of squalamine and their utility in inhibiting neovascularization and endothelial cell proliferation.

CROSS REFERENCE TO RELATED APPLICATIONS [0001] This application claims the benefit of U.S. provisional applications 60/942,364 filed on Jun. 6, 2007, and 60/944,696 filed on Jun. 18, 2007, which are incorporated by reference as if fully set forth. TECHNOLOGY FIELD [0002] The subject matter disclosed is related to wireless communications. More particularly, the subject matter is related to supporting media independent handover (MIH). BACKGROUND [0003] The IEEE 802.21 standard provides a uniform set of functionalities that help enable and enhance handovers across different link layer technologies. IEEE 802.21 defines three main services available to Mobility Management applications, such as Client Mobile Internet Protocol (Client MIP) or Proxy MIP. Referring to FIG. 1 , these services are the Event Service 100 , the Information Service 105 and the Command Service 110 . These services aid in the management of handover operations, system discovery and system selection by providing information and triggers from lower layers 115 to upper layers 120 via a media independent handover (MIH) function (MIHF) 125 . [0004] At a high level, this involves an upper layer MIH User which can communicate with an MIH Function 125 (either locally or remotely over some transport medium) through link-independent Event Service 100 , Information Service 105 and Command Service 110 . The MIH Function 125 , in turn, will interact with link-layer devices through the use technology-specific primitives; the functionalities expected from these technology-specific primitives are defined in the 802.21 standard. While FIG. 1 shows MIHF 125 as a middle layer in a protocol stack, MIHF 125 may also be implemented as an MIH plane that is capable of exchanging information and triggers directly with different layers of the protocol stack. [0005] The Third Generation Partnership Project (3GPP) has identified three principles that describe how inter-system handovers between 3GPP and non-3GPP access (e.g. 3GPP2, IEEE 802.11 WLAN, IEEE 802.16 WiMAX, etc.) should be handled. However, these principles do not address how two different accesses can be integrated in order to allow handover execution. The first principle applies in multiple RAT scenarios where the wireless transmit/receive unit (WTRU) uses a single radio access technology (RAT) for all in-progress services. The second principle is that the Inter-RAT handover decision is made and the handover command is sent by the serving Radio Access Network (RAN). The target RAN may exercise admission control to the WTRUs that are handed over. The third principle is that the serving RAN receives information from the target RAN that can be included in the handover command. [0006] All these principles can be met by using the handover (HO) service provided by the 802.21 standard. This is especially needed when handover commands requesting a switch over toward or from a 3GPP based access is required, for example, when a handover takes place between IEEE 802.16 or WiMAX accesses and 3GPP accesses, or between IEEE 802.11 or WLAN systems and 3GPP systems. [0007] FIG. 2 depicts a typical GSM Edge Radio Access Network—UMTS Terrestrial Radio Access Network (GERAN-UTRAN) 3GPP packet switched (PS-domain) Inter-RAT architecture 200 . Referring to FIG. 2 , the source network includes a serving GPRS support node SGSN 205 , a base station controller/radio network controller (BSC/RNC) 210 , and a base transceiver station (BTS)/Node B 215 . The BSC/RNC 210 communicates with the SGSN 205 through a Gb/IuPS interface 220 . In addition, the BSC/RNC 210 communicates with the BTS/Node B 215 through an Abis/Iub interface 225 . The target network includes a SGSN 230 , a BSC/RNC 235 , and a BTS/Node B 240 . The BSC/RNC 235 communicates with the SGSN 230 through a Gb/IuPS interface 245 . The BSC/RNC 235 communicates with the BTS/Node B 240 through an Abis/Iub interface 250 . The source and target SGSNs 205 , 230 communicate through a Gn interface 255 . [0008] Referring to FIG. 2 , it is the source BSC/RNC 210 that controls the handover. The mobile node (MN) 260 is requested to take measurements in the target network and, upon meeting the handover conditions, the source BSC/RNC 210 requests the target BSC/RNC 235 to prepare the resources for the MN 260 . The target BSC/RNC 235 performs admission control and responds with the new resource allocation. Once the new resources have been allocated, the source BSC/RNC 210 commands the MN 260 to handover to the new network. Upon detecting the MN 260 in the new network, the target BSC/RNC 235 informs the source BSC/RNC 210 of the handover completion. [0009] In order to perform heterogeneous handover between a 3GPP and non-3GPP network, the network architecture must provide capability for an MIH User to acquire measurement reports and capability for an MIH Function to reserve link layer resources through the use of standardized MIH primitives and messages. While the 802.21 standard provides mechanisms to obtain such measurement reports, query for resources, reserve these resources, execute the handover and inform the peer network about the completion of the action, the mechanisms have deficiencies that deprive implementers from the use of key functionalities and from complete control of the measurement-reporting process. This is specifically true for handover between 3GPP (e.g. GERAN, UTRAN and LTE) and non-3GPP networks, which are also known as Inter-Radio Access Technology (Inter-RAT) handovers. [0010] When two peer networks are to perform a handover, typically based on Mobile Node (MN) (also referred to as User Equipment or UE) measurement reports, the network instructs the MN to switch to another cell and indicates what configuration to use in the new cell. This implies that the Inter-RAT handover decision is made by the serving Radio Access Network (RAN), whereas the target RAN may exercise admission control on the MN that is being handed over. [0011] Hence, the sequence of events is 1) a Query phase used to determine the status of resources at both source and target networks before taking a handover decision, 2) a Preparation phase where resources are reserved at the target network once a handover decision has been taken, 3) an Execution phase when the handover commands are sent and performed, and 4) a Completion phase when the result of the handover is informed and the original resources are released. [0012] The IEEE 802.21 specification defines messages that can be used to perform the actions described above. However, the functionality provided by the currently defined messages is insufficient to convey all the required information between source and target networks, especially in the case of 3GPP to non-3GPP handover (and vice versa). It would therefore be desirable to provide messages to convey all the required information between source and target networks without compromising functionality. In order to perform heterogeneous handover between a 3GPP and non-3GPP network, it would also be desirable to design a network architecture to provide capability for an MIH User to acquire measurement reports and capability for an MIH Function to reserve link layer resources through the use of standardized MIH primitives and messages. SUMMARY [0013] A method and apparatus for access-independent mobility management. The method and apparatus are used in handover between 3GPP and non-3GPP networks which use enhanced media independent handover functionalities. BRIEF DESCRIPTION OF THE DRAWINGS [0014] A more detailed understanding may be had from the following description, given by way of example and to be understood in conjunction with the accompanying drawings wherein: [0015] FIG. 1 is an IEEE 802.21 protocol architecture according to the prior art; [0016] FIG. 2 is a block diagram for a 3GPP PS-domain Inter-RAT architecture according to the prior art; [0017] FIG. 3 is a block diagram of a system performing Inter-RAT Handover with a Proxy MIH Node; [0018] FIG. 4 is a block diagram of a system performing for an Inter-RAT Handover with an MIH-capable SGSN/Network Controller; [0019] FIG. 5 is a block diagram of a system performing an Inter-RAT Handover with MIH Server; [0020] FIG. 6 is a block diagram of a system performing Inter-RAT Handover; [0021] FIG. 7 is a block diagram of a system using media independent normalizing functions to interpret 3GPP commands and map their functionality into equivalent generic handover commands; [0022] FIG. 8 is a block diagram of a system using media independent normalizing functions to interpret 3GPP commands and map their functionality into equivalent generic handover commands; [0023] FIG. 9 shows a block diagram of a roaming scenario where the MN is in a visited network; [0024] FIG. 10 is a WTRU, Access Point (AP) or Point of Access (PoA) and a Point of Service (PoS) or MIH Server configured to perform heterogeneous handover between a 3GPP and non-3GPP network using MIH messaging; [0025] FIG. 11 is a signal diagram of a system performing Inter-RAT Handover using media independent normalizing functions; [0026] FIG. 12 is a signal diagram of a system performing Inter-RAT Handover using media independent normalizing functions and single-radio with on-off techniques; and [0027] FIG. 13 is a signal diagram of a system performing Inter-RAT Handover using media independent normalizing functions and multi-radio techniques. DETAILED DESCRIPTION [0028] When referred to hereafter, the terminology “wireless transmit/receive unit (WTRU)” includes but is not limited to a user equipment (UE), a mobile station, a fixed or mobile subscriber unit, a pager, a cellular telephone, a personal digital assistant (PDA), a computer, a mobile node (MN), or any other type of user device capable of operating in a wireless environment. When referred to hereafter, the terminology “base station” includes but is not limited to a Node-B, an Enhanced Node-B (eNB), a site controller, an access point (AP), or any other type of interfacing device capable of operating in a wireless environment. [0029] The embodiments below are described in reference to the 802.21 protocol and messages for simplicity. Although the embodiments described below refer to messages defined in the 802.21 protocol, the concepts can be applied messages defined in other technologies containing similar information elements to 802.21 messages. [0030] IEEE 802.21 services, for example, and in particular Command and Information services, can be used to integrate multiple access technologies. This includes system architecture that show where the Media Independent Handover function can be placed in order to allow this integration. Also included is a mechanism that shows how mobility principles, as outlined by 3GPP standards, can be met using the proposed architecture. Through the use of services provided by the MIH Function, a mobility mechanism supporting Handover between 3GPP and non-3GPP access can be realized. The location of the MIH function within the 3GPP architecture is logically distributed and it might depend on the level of integration that is desired, that is, whether a tight coupling or a loose coupling scenario is being addressed. [0031] Three logical components, i.e., the MME, the Gateway, and the IP server, can be identified. These logical components can communicate amongst each other or act independently depending on specific deployment scenarios. Logically the MIH function could also reside within a specific access if a particular deployment warrants it. [0032] The basic functionality for the 3GPP architecture is defined in FIG. 2 above. Using the basic architecture from FIG. 2 , the following three network architectures can be derived for the non-3GPP case to support heterogeneous handover. [0033] FIG. 3 shows one possible architecture 300 that can be used to support the heterogeneous handover between 3GPP and non-3GPP networks. Referring to FIG. 3 , the source network includes a SGSN 305 , a base station controller/radio network controller (BSC/RNC) 310 , and a BTS/Node B 315 . The BSC/RNC 310 communicates with the SGSN 305 through a Gb/IuPS interface 320 . In addition, the BSC/RNC 310 communicates with the BTS/Node B 315 through an Abis/Iub interface 325 . The target network includes a Generic Network Gateway 330 , a Generic Network Controller 335 , and a Generic Base Station 340 . [0034] Referring to FIG. 3 , an 802.21 MIH node 345 is used to translate and act as a proxy between the Generic Network Gateway 330 and the 3GPP SGSN 305 . If a conventional SGSN is used, the handover messages communicated between the MIH Proxy 345 and the SGSN 305 would be the same as described in the 3GPP Gn interface 350 . If the network is small, or the SGSN 305 and BSC/RNC 310 are collocated, the MIH Proxy 345 could connect directly to the BSC/RNC 310 by using Iu messages 355 . [0035] FIG. 4 shows another possible network architecture 400 to perform an Inter-RAT Handover with MIH-capable SGSN/Network Controller. Referring to FIG. 4 , it is assumed that the SGSN 410 and Generic Network Gateway 420 implement MIH capabilities 415 , 425 , and therefore are capable of communicating one to another with MIH messages 430 , such as messages defined in the 802.21 protocol or messages defined in other technologies containing similar information elements to 802.21 messages. [0036] A similar approach could be applied where the Generic Network Controller 435 and BSC/RNC 440 were MIH-capable. For this approach, these two nodes would be able to communicate with MIH messages without passing through the gateways. For simplicity, this approach is not shown in FIG. 4 . [0037] FIG. 5 shows an alternative network architecture 500 for an Inter-RAT Handover with MIH Server. In this architecture, the MIH Server 510 acts on behalf of the Network Controller for taking handover decisions (e.g. as source Network Controller) and setting up the resources at the target network. In this figure, it is shown that the MIH Server 510 can communicate to the SGSN 515 , for example, through a Gn interface 520 , and/or to the BSC/RNC 525 , for example, through a Gb/Iu interface 530 . Also in this figure, the MIH Server 510 communicates directly to the mobile Node (MN) 535 via L2/L3 protocols (e.g. 802.11, 802.16, IP, etc.) 540 . [0038] In order to support heterogeneous handover between a 3GPP and a non-3GPP network, media independent handover messages can be used. For instance, the existing 802.21 standard messages or other technologies standards can be updated to include the following messages: MIH_N2N_HO_Commit request; and MIH_N2N_HO_Commit response. [0041] By including these two messages, the MIH network functionality (or similar network functionality) has the capability to reserve resources when two networks control the handover, similar to the 3GPP networks. [0042] Although the 802.21 standard, for example, can be updated to include the required messages, the contents of these messages do not fulfill the requirements of the 3GPP network handover. Hence, an enhancement to the MIH messages is required to support handovers between 3GPP and non-3GPP networks. This enhancement would follow the Inter-RAT Handover (GERAN/UTRAN) philosophy described in the background section above. [0043] The enhanced messages and their encoding, e.g., TLV IEs (Type-Length-Value Information Elements), are discussed in the embodiments below. Where the networks are not pre-configured with each other's parameters, the source network can request the target network about the available resources (e.g. cell list, cell parameters, etc.). For this, the source network can either ask the target to report on all available resources, or on a specific type of network. [0044] In one embodiment, this information could be included in the following MIH Messages: N2N Query Resources Request (from source to target network to request reporting on available resources that could be used by the source to handover). MN HO Query Request (from mobile to target network to request reporting on available resources that could be used by the source to handover). [0047] One possibility is to use the Network Type element to request information about a specific network. Another possibility is to include the network information as part of the Available Resource field of the above mentioned message as a suggestion from the source. [0048] FIG. 6 shows how the updated handover messages can be used to perform an Inter-RAT Handover 600 . Before the handover process starts, it is required for the WTRU to start searching neighboring cells 605 and provide measurements. In order to perform such measurements for 3GPP GERAN/UTRAN/LTE or non-3GPP networks, neighbor list and measurement information is required by WTRU to take measurements on neighbor cells. Thresholds and event criteria (i.e., when to report measurements), periodicity of measurements, and number of cells to report can optionally be included in this information. [0049] In one embodiment, the information required for 3GPP GERAN/UTRAN/LTE or non-3GPP networks could be included in the following enhanced MIH Messages: N2N Query Resources Response (from target to source network to inform the available cells that should be scanned in the network); Net HO Query Request (from source network controller to MN to let the MN know which cells to monitor); and MIH Scan Request (from source network controller to MN to let the MN know which cells to monitor). [0053] One possibility is to include the information as part of the Available Resource field of the above mentioned enhanced messages. [0054] Referring to FIG. 6 , when MIH server requests reports 610 , or the WTRU independently triggers a measurement report 605 , the required information, such as the cell ID of the best cell or list of best cells, could be included in the following enhanced MIH Messages: Link Parameter Report (from MN to the network to report on measurements); Net HO Query Response (from MN to the network to respond to the query request and report on measurements); and MIH Scan Response (from MN to the network to respond to the request and report on measurements). [0058] One possibility is to include the information as part of the Link Parameters, Link Resource, or Scan Response fields of the above mentioned enhanced messages. [0059] Referring to FIG. 6 , upon receiving a measurement report, the MIH server performs reservation of resources for the target cell 615 . To perform a reservation, the MIH can communicate directly to the target SGSN or mobility management entity (MME) 620 or alternatively to the eNB, RNC or MSC 625 by making use of existing handover messages, such as “Prepare Handover”. [0060] The required information to reserve resources on the target network could be included in the following enhanced MIH Messages: N2N HO Commit Request (from source to target network to request reservation of the resources); and MN HO Commit Request (from MN to network to request reservation of the resources). [0063] This information could, in one embodiment, be included in the Query Resource, or Reserve Resource fields of the above mentioned enhanced messages. [0064] Referring to FIG. 6 , once the resources have been reserved by the target network 630 , the source network (or WTRU) is informed about the successful reservation of resources 635 so that the handover can take place. Hence, the information required by the MN to make the connection to the new network could be included in the following updated MIH Messages: N2N HO Commit Response (from target to source network to report reservation of the resources); MN HO Commit Response (from network to MN to report reservation of the resources); and Net HO Commit Request (from the network to the MN to report reservation of resources and command the MN to handover to these resources). [0068] This information could, in one embodiment, be included in the Query Resource, or Reserve Resource fields of the above mentioned messages. [0069] Referring to FIG. 6 , once the reservation of resources is complete, the handover information is sent to the MN or WTRU 640 in order to perform handover to the target network 645 . Once the handover is complete, handover complete messages can be sent 650 to re-route traffic through the new network and release resources from the source network. [0070] Depending on the type of network, the handover can be performed in a variety of ways. For GSM, once the BSC has reserved the radio resources of GERAN cell resources it has to give the necessary information for the WTRU to complete the handover and synchronize to the new cell. This information is transmitted to the WTRU via the source network in a transparent container. Such type of transparent container can be used in other types of network to convey the information of the radio resources either from source to target or vice versa. [0071] The following information for the WTRU, transmitted in a transparent container, could be contained in the MIH message: Synchronization Indication (SI); Normal Cell Indication (NCI); ARFCN, BSIC-BCCH frequency and BSIC of new cell; CCN Support Description; Frequency parameters; Extended dynamic allocation; Network Control Order; RLC reset; Packet timing Advance; UL control timeslot; GPRS, EGPRS mode; and UL/DL TBFs (PFI, TFI assignment, TBF timeslot allocation, RLC mode, USF allocation); Optional: NAS container. [0086] For UTRAN, once the RNC has reserved the radio resources for the cell id, it has to give the necessary information for the Mobile station to complete the handover and synchronize to the new cell. This information is transmitted to the WTRU via the source network in a transparent container. [0087] The following information, transmitted in a transparent container in a MIH message, is required by the WTRU to make the connection to the 3G cell: WTRU identities (U-RNTI, H-RNTI, E-RNTI); Ciphering algorithm; RB information elements (SRB information to setup list, RAB information to setup list); UL/DL transport channel information (UL/DL Transport channel information common for all transport channels, Added or Reconfigured TrCH information UL/DL); UL radio resources (Uplink DPCH info, E-DCH Info); DL radio resources (Downlink HS-PDSCH Information, Downlink information per radio link, Downlink information common for all radio links); Frequency info; and Maximum allowed UL tx power. [0096] In addition, the RNC may provide information for Commit time/activation for synchronous handovers. [0097] Alternatively, predefined configurations can be used if the WTRU supports them. A predefined configuration will require less information to be transmitted to the WTRU: Default configuration mode (FDD, TDD); Default configuration identity; RAB info; and UL DPCH info. [0102] The RNC may also provide MIH Complete Request/Response Messages. Once the MN has been handed over from the source to the target network, handover complete messages are sent to re-route traffic through the new network and release resources from the source network. [0103] In one embodiment, this information could be included in the following enhanced MIH Messages: Net HO Commit Response; N2N Complete Request; and N2N Complete Response. [0107] For LTE and other 3GPP technologies such as WCDMA and GERAN media independent normalizing functions can be used to interpret 3GPP commands and map their functionality into equivalent generic handover commands, such as the ones described in IEEE 802.21. FIGS. 7 , 8 and 9 show how this media independent handover function can be logically placed, for example, within the PDN Gateway 710 as this is the central point of contact across multiple access systems. The 3GPP network 715 shown in FIG. 7 includes an MME 720 capable of supporting E-UTRAN 720 communications. The MME 720 is also in communication with a 2G/3G SGSN 725 , which is capable of supporting UTRAN 730 and GERAN 735 communications. The non-3GPP network 740 includes an ePDG 745 capable of supporting untrusted non-3GPP access 750 . The trusted non-3GPP access 755 is in direct communication with the PDN Gateway 710 . [0108] As described in FIG. 8 , the WTRU 805 remains under the domain of 3GPP handover mechanism while the current connection is progress. The Target MIH PoS PDN Gateway 810 serves as the central point of contact between the 3GPP 815 and non-3GPP networks 820 . The source SGSN/MME 825 can use Forward_Relocation_Req 830 and Forward_Relocation_Complete 835 messages to communicate with the Target MIH PoS PDN Gateway 810 . The Trusted Non-3GPP Access 840 can use MIH_N2N_HO_CandidateQuery/MIH_N2N_HO_Commit request 845 and MIH_N2N_HO_Candidate_Query/MIH_N2N_HO_Commit response 850 messages to communicate with the Target MIH PoS PDN Gateway 810 . Similarly, the MN can use HO Commit and Query request and response types of messages to trigger or initiate the handover and to obtain the required information for handover once the preparation is finished. [0109] FIG. 9 shows an example of a roaming scenario 900 where the MN 905 is in a visited network 910 . In this scenario, there are two gateways in which the MIH 915 could reside, the Serving Gateway 920 and the Anchor Gateway 930 . This scenario may also include an IP server 940 which can communicate with the MN 905 , for example using an IP interface. The MIH functionality 915 may also be located in the MME 950 . This example is also be applicable to the home scenario. In an alternative embodiment, the MIH 915 may be located in E-UTRAN 960 . [0110] The WTRU may or may not be able to simultaneously support multi radio capabilities or only one radio technology at time. If multiple radio capabilities are supported either by using multi-radio or single-radio with on-off techniques, the WTRU might be able to measure radio environments from multiple accesses while still connected to the current access. Normalized measurement reporting capabilities, such as the ones described in 802.21, could be used to provide a service access point for measurement collection purposes, exposing a unified interface regardless of the underlying technology. [0111] The WTRU might also rely on information provided via higher layers over the current access by using information services such as the ones provided by IEEE 802.21. This information allows the WTRU to request access relocation, even when no specific measurements are provided. [0112] When preparing and reserving radio resources, the MIH Function is able to map the relocation request to a suitable MIH Command. This allows the target access system to exercise admission control functions prior to granting resources. The command that triggers the handover from the 3GPP access is generated entirely according to 3GPP specifications, possibly using information provided by the target access system via MIH mapping. [0113] Table 1 below shows a possible mapping between the MIH, e.g., enhanced 802.21, and 3GPP GERAN/UTRAN/LTE messages that could be used, for instance, by the proxy function. [0000] TABLE 1 Air LTE 802.21 Gn Iu Gb Interface (S11/S3/S4) N2N Commit Forward Relocation PS Forward Request Relocation Request Handover Relocation Request Required Request N2N Commit Forward Relocation PS Forward Response Relocation Request Handover Relocation Response Ack Required Response Ack Net HO Commit PS PS HO Request Handover Command Command Net HO Commit HO to HO to E- Response UTRAN UTRAN Complete Complete N2N HO Complete Forward Relocation Forward Request Relocation Complete Relocation Complete Complete N2N HO Complete Forward Forward Response Relocation Relocation Complete Complete Ack ACK N2N Commit Update Request Bearer Request N2N Commit Update Response Bearer Response N2N Commit Forward Request SRNS Context N2N Commit Forward Response SRNS Context ACK N2N_HO_Candidate — Forward Query Request Relocation Request N2N_HO_Candidate — Forward Query Relocation Response Response N2N HO Complete Update Request Bearer Request N2N HO Complete Update Response Bearer Response [0114] Tables 2-5 below show a possible realization combination of the message encoding that would carry the above mentioned parameters in a type-length-value (TLV) format. [0000] TABLE 2 System Parameters List Type Length Value XXX Variable Structure consisting of 1) Network Type, and 2) Network Specific System Parameters [0000] TABLE 3 Network Type Type Length Value XXX 8 Network Type and Revision as defined in 802.21 standard [0000] TABLE 4 Network Specific Parameters Type Length Value XXX Variable Network Specific System Parameters. 802.16: UCD, DCD, UIUC, DIUC GSM/GRPS/EDGE (GERAN): (defined depending on message type) 3GPP (UTRAN): (defined depending on message type) [0000] TABLE 5 HANDOVER COMPLETION Type Length Value XXX Variable Parameter Integer Type of network type 0: IEEE 802.16 1: GERAN 2: 3GPP 3-7: Reserved System Variable Depending on the parameter type parameter 0: UCD, DCD, UIUC, DIUC Value 1: (defined depending on message type) 2: (defined depending on message type) 3-255: Reserved [0115] FIG. 10 is a WTRU 1000 and access point 1005 configured to implement the IEEE 802.21 Inter-RAT Handover as described above. WTRU 1000 includes a processor 1010 , an MIH function 1015 , and a plurality of transceivers 1020 a . . . 1020 n , each configured to operate using a different radio access technology and protocol. The processor 1010 and MIH function 1015 are configured to operate protocol stacks according to the above described embodiments. Further, the Processor 1010 and MIH function 1015 are capable of generating enhanced messages as described above, for example, with reference to FIG. 8 . The processor 1010 and MIH function 1015 are further configured to implement IEEE 802.21 protocols for MIH peer messaging. The IEEE 802.21 messages may be transmitted to MIH peers via any of the plurality of transceivers 1020 a . . . 1020 n . The processor 1010 and MIH function 1015 are further configured to implement local IEEE 802.21, for example for the IEEE 802.21 Command service. The transformation of MIH messages, and the extraction of MIH messages from received messages may be performed by either processor 1010 or MIH function 1015 , or by a combination of the two. [0116] Access point 1005 includes a processor 1025 , an MIH function 1030 , and a transceiver 1035 . The access point 1005 communicates with WTRU 1000 via air interface 1040 . The processor 1025 of the access point 1005 processes received IEEE 802.21 messages received from WTRU 1000 via transceiver 1035 . The processor 1025 and MIH function 1030 of the access point 1005 are further capable of generating enhanced messages as described above, for example, with reference to FIG. 8 . The processor 1025 and MIH function 1030 are further configured to implement IEEE 802.21 protocols for MIH peer messaging, such as messaging between the access point 1005 and an MIH server (MIHS) 1045 , or a PoS (not shown). The transformation of MIH message, and the extraction of MIH messages from received messages may be performed by either processor 1025 or MIH function 1030 , or by a combination of the two. [0117] FIG. 11 is a signal diagram of a system 1100 performing Inter-RAT Handover using 802.21 media independent normalizing functions. The system includes a WTRU 1110 , a source network 1020 , an MIH Proxy 1130 and a target network 1140 . [0118] Referring to FIG. 11 , the WTRU 1110 searches neighboring cells 1115 and provides a measurement report 1125 to the MIH Proxy 1130 via the source network 1120 . The MIH Proxy 1130 performs reservation of resources 1135 for the target network 1140 . Once the resources are reserved 1150 in the target network 1140 , the source network 1120 is informed of the successful reservation of resources 1155 via the MIH Proxy 1130 . The handover information 1160 is then sent from the source network 1120 to the WTRU 1110 . The WTRU 1110 then performs the handover 1165 to the target network 1140 . The target network 1140 then sends a handover complete message 1170 to the source network 1120 . [0119] FIG. 12 is a signal diagram of a system 1200 performing Inter-RAT Handover using 802.21 media independent normalizing functions and single-radio with on-off techniques. The system includes a WTRU 1210 , a source network 1220 , an MIH server 1230 , and a target network 1240 . [0120] Referring to FIG. 12 , the WTRU 1210 searches neighboring cells 1215 and provides neighbor information 1225 to the MIH. Optionally, the WTRU 1210 may be triggered by an MIH request 1235 to begin searching neighboring cells 1215 . Upon receiving the neighbor information 1225 , the MIH server 1230 performs reservation of resources for the target cell 1245 via the source network 1220 . Once the resources are reserved 1250 in the target network 1240 , the source network 1220 is informed of the successful reservation of resources 1255 . The handover information 1260 is then sent from the source network 1220 to the WTRU 1210 . The WTRU 1210 then performs the handover 1265 to the target network 1240 . The target network 1240 then sends a handover complete message 1270 to the source network 1220 . [0121] FIG. 13 is a signal diagram of a system 1300 performing Inter-RAT Handover using 802.21 media independent normalizing functions and multi-radio techniques. The system includes a WTRU 1310 , a source network 1320 , an MIH server 1330 , and a target network 1340 . [0122] Referring to FIG. 13 , the WTRU 1310 searches neighboring cells 1315 and provides neighbor information 1325 to the MIH. Optionally, the WTRU 1310 may be triggered by an MIH request 1335 to begin searching neighboring cells 1315 . Upon receiving the neighbor information 1325 , the MIH server 1330 performs reservation of resources for the target cell 1345 . Once the resources are reserved 1350 in the target network 1340 , the source network 1320 is informed of the successful reservation of resources 1355 . The handover information 1360 is then sent from the source network 1320 to the WTRU 1310 . The WTRU 1310 then performs the handover 1365 to the target network 1340 . The target network 1340 then sends a handover complete message 1370 to the source network 1320 . [0123] Note that the target network 1340 can also send the resource reservation directly to the WTRU 1310 using the target network air interface (not shown), without having to go through the source network 1320 . The WTRU 1310 has dual radio so it can receive from the target network 1340 without service interruption from the source network 1320 . The source network 1320 should be notified that the handover has been completed, but either the target network 1340 or the WTRU 1310 can release the connection. In this situation, the MIH server 1330 informs the WTRU 1310 to perform the handover based either on dynamic measurements or static policies. The WTRU 1310 then proceeds to reserve and connect directly to the target network 1340 without passing through the MIH server 1330 or the source network 1320 . [0124] Although the features and elements of the present invention are described in the preferred embodiments in particular combinations, each feature or element can be used alone without the other features and elements of the preferred embodiments or in various combinations with or without other features and elements of the present invention. The methods or flow charts provided in the present invention may be implemented in a computer program, software, or firmware tangibly embodied in a computer-readable storage medium for execution by a general purpose computer or a processor. Examples of computer-readable storage mediums include a read only memory (ROM), a random access memory (RAM), a register, cache memory, semiconductor memory devices, magnetic media such as internal hard disks and removable disks, magneto-optical media, and optical media such as CD-ROM disks, and digital versatile disks (DVDs). [0125] Suitable processors include, by way of example, a general purpose processor, a special purpose processor, a conventional processor, a digital signal processor (DSP), a plurality of microprocessors, one or more microprocessors in association with a DSP core, a controller, a microcontroller, Application Specific Integrated Circuits (ASICs), Field Programmable Gate Arrays (FPGAs) circuits, any other type of integrated circuit (IC), and/or a state machine. [0126] A processor in association with software may be used to implement a radio frequency transceiver for use in a wireless transmit receive unit (WTRU), user equipment (UE), terminal, base station, radio network controller (RNC), or any host computer. The WTRU may be used in conjunction with modules, implemented in hardware and/or software, such as a camera, a video camera module, a videophone, a speakerphone, a vibration device, a speaker, a microphone, a television transceiver, a hands free headset, a keyboard, a Bluetooth® module, a frequency modulated (FM) radio unit, a liquid crystal display (LCD) display unit, an organic light-emitting diode (OLED) display unit, a digital music player, a media player, a video game player module, an Internet browser, and/or any wireless local area network (WLAN) module.

Methods and mechanisms enhance heterogeneous media independent handover (MIH) between different link layer technologies. Embodiments include using an MIH proxy entity, MIH capable network controller, and an MIH server. Enhancements are made to the query phase, preparation phase, execution phase and completion phase by including required information in MIH messages.

FIELD OF INVENTION The invention relates generally to the art of digital communications and more specifically to a system for minimizing the average latency in transporting messages, such as packets or frames, which are segmented into a plurality of smaller cells for transport across a network. BACKGROUND OF INVENTION Asynchronous transfer mode (hereinafter “ATM”) service inter-networking protocols enable data or messages formatted according to a non-ATM data communication protocol to be transported across an ATM network. For example, the Frame Relay Forum FRF.5 protocol specifies how a relatively large, variable length, frame relay packet should be segmented into a plurality of ATM-like, fixed-size, cells for transport across an ATM network. Such protocols necessarily define how the ATM Adaption Layer (AAL) should be provisioned since this layer of the ATM/B-ISDN protocol stack, as defined by ITU Recommendation I.321 and shown in FIG. 1, is responsible for adapting the services provided by the ATM Layer, which provides basic ATM cell transport functions, to higher layers, e.g. frame relay bearer service. FIG. 2 illustrates a generic version of the AAL in greater detail. As shown in FIG. 2, some versions of the AAL, such as AAL 3 / 4 and AAL 5 , include a convergence sublayer (CS) and a segmentation and reassembly sublayer (SAR). The CS, which sits directly above the SAR and below the AAL Service Access Point (SAP), aligns the SDUs and adds overhead information. The CS may also provide service specific signalling or data link functions. The SAR, when operating in a message mode, segments a single AAL SAR Service Data Unit (hereinafter “SDU”), such as a variable length frame packet, into a plurality of AAL SAR Protocol Data Units (hereinafter “PDU”), each of which essentially forms the payload of an ATM cell transmitted across an ATM network. Conversely, at the destination, the destination SAR requires that all of the PDUs composing an SDU be passed from the ATM Layer to the destination SAR before it can reassemble the SDU and, ignoring the role of the convergence sublayer, indicate reception of the SDU to the higher layer using the AAL. Thus, the latency in transmitting an SDU from a first point to a second point in a network can be defined as the time from which the transmission of the SDU is first requested until the time the last PDU arrives at the destination SAR and the SDU is reassembled. In other words, latency can be defined as the time required to transmit the SDU from an originating AAL SAP to a destination AAL SAP. This latency is entirely characterized by the amount of time required to propagate the last PDU of an SDU across the ATM network—the time required to propagate any other PDU before the last PDU of an SDU is of no consequence at the destination AAL SAP. Latency manifests itself as sluggishness or slow response time in interactive-type communications. For example, if one were sending joystick instructions across a network during the course of an interactive game played there over, a long latency would, in the absence of other aggravating factors, result in a noticeable time period between the physical movement of the joystick and the corresponding computer action. Accordingly, it is desirable to minimize latency for interactive telecommunications applications. Latency is affected by the service discipline used to schedule or multiplex PDUs corresponding to SDUs from a plurality of virtual connections (VCs) into a single cell stream for transmission across the Physical Layer (PHY) of the ATM network. FIG. 3 shows how an ATM Layer 11 provides a SAP 10 to each of several VCs, each of which has its own AAL 12 (i.e., the AAL is invoked in parallel instances). The ATM Layer 11 , in turn, uses a single SAP 14 into a PHY 16 . One role of the ATM Layer 11 is to accept requests of PDUs 17 from each SAP 10 and to multiplex these PDUs into a single cell stream 18 such that the timing of the transmission of each of the PDUs conforms to predetermined traffic parameters assigned to its respective VC. FIG. 3 illustrates a condition where each VC generates a burst 20 of several ATM PDU requests at the ATM SAP 10 , wherein each such burst corresponds to a single SDU 22 , such that there is an overlap in the transmission periods of the SAR SDUs from the ATM Layer 11 to the PHY 16 . The ATM PDUs 17 received from the ATM SAP 10 must therefore be queued, and then the ATM PDUs from each of the different ATM SAPs 10 must be multiplexed in some order onto the single stream 18 of ATM PDUs passed to the PHY SAP 14 . Given this set of PDUs which have been requested over several ATM Layer SAPs 10 , and subject to the constraints of satisfying the traffic parameters of each VC, it is often desired to minimize the average amount of latency experienced per unit of SDU data (i.e., per PDU) for various types of ATM service categories. As shown in FIG. 3, a typical ATM Layer implementation might use round-robin ordering in sending the PDUs 17 to the PHY SAP 14 from each ATM SAP 10 . This would result in each corresponding SDU 22 using an equal fraction of the PHY bandwidth while the PDUs for each SDU are being transmitted. This is shown, for instance, in the bandwidth occupancy chart of FIG. 4A for the situation where two VCs each request a burst of the same number of PDUs at about the same time, wherein each VC has a PCR equal to 100% of the available bandwidth. (A “bandwidth occupancy chart” is a chart with time on the horizontal axis, and bandwidth on the vertical axis. Each SDU sent on an ATM virtual connection is shown as a shaded region on such a chart. The net height of the region at a particular time shows the amount of bandwidth occupied by the transmission of the SDU at that time; the leftmost and rightmost extent of the region gives the time at which the first and last PDUs for the SDU are transmitted, respectively; and the total area of the region gives the size of the SDU. Unshaded regions in these charts represent the proportion of unused PHY bandwidth, for which the ATM Layer will be sending idle cells.) This ordering is not optimal with respect to the average amount of latency experienced per unit of SDU data. SUMMARY OF INVENTION Broadly speaking, the invention seeks to minimize or reduce the average per unit latency in transporting messages which are decomposed into a plurality of smaller data units for transport across a network. One aspect of the invention relates to an apparatus for transmitting messages associated with a plurality of variable bit rate connections, each of which is associated with a traffic contract which defines compliance thereto as conformance to a leaky bucket algorithm. The apparatus comprises transmission equipment for receiving multiple messages, segmenting each received message into one or more data units, and multiplexing such data units from various connections into a single stream for transport over a physical interface at an output transmission rate. A bandwidth allocation means is associated with the transmission equipment for dynamically allocating a portion of the output transmission rate to any connection. A scheduler is connected to the bandwidth allocation means for scheduling the transfer of messages to the single stream of the transmission equipment and for allocating a portion of the output transmission rate to each connection at the time its message is transferred to the transmission equipment. The portion of the output transmission rate for a given connection is substantially equal to 1/T, T being computed as T ← max  ( T S + X - τ S N - 1 , T p , T L ) , if        N > 1 and T← max( T p , T L ), if N =1 where T S is a period corresponding to a constant sustained transmission rate, T p is a period corresponding to a peak transmission rate, τ S is a burst tolerance, N is the number of data units in the message, X is a fill level of the leaky bucket associated with the given connection, and T L corresponds to all unused bandwidth allocated to a service class to which the given connection belongs. Another aspect of the invention relates to an apparatus for transmitting messages associated with a plurality of variable bit rate connections, each of which is associated with a traffic contract which defines compliance thereto as conformance to a leaky bucket algorithm. The apparatus comprises transmission equipment for receiving multiple messages, the transmission equipment including a segmentation means for segmenting each received message into one or more data units and a multiplexing means for multiplexing such data units from various connections into a single stream for transport over a physical interface at an output transmission rate. A bandwidth allocation means is associated with the transmission equipment for dynamically allocating a portion of the output transmission rate to any connection. A scheduler, connected to the bandwidth allocation means, is provided for scheduling the transfer of messages to the single stream of the transmission equipment. The scheduler preferably transfers messages to the single stream of the transmission equipment by receiving the messages from higher layer networking services and transferring the messages to the segmentation means which substantially immediately transfers the messages to the multiplexing means. Alternatively, the scheduler may transfer messages to the single stream of the transmission equipment by receiving the data units from the segmentation means and transferring groups of data units to the multiplexing means, wherein each such group corresponds to a message. The scheduler is enabled to transfer the messages according to an order corresponding to the level of a leaky bucket associated with each corresponding connection in order to reduce the interleaving of data units from various connections over the single stream. In the preferred embodiment, the transmission equipment comprises a segmentation means and a multiplexing means for respectively segmenting each received message into data units and multiplexing the data units into the aforesaid single stream. The scheduler preferably transfers messages to the single stream of the transmission equipment by receiving the messages from higher layer networking services and transferring the messages to the segmentation means which substantially immediately transfers the messages to the multiplexing means. Alternatively, the scheduler may transfer messages to the single stream of the transmission equipment by receiving the data units from the segmentation means and transferring groups of data units to the multiplexing means, wherein each such group corresponds to a message. In the preferred embodiment, the transmission equipment includes an indication means for informing the scheduler when a message has been fully transmitted by the transmission equipment. When indication means is actuated, the scheduler is enabled to set a new transmission rate for a given connection during the time a message associated therewith is in the process of being transmitted by the transmission equipment. The new transmission rate is substantially equal to 1/T, T being computed as T ← max  ( T S + X est - τ S R est - 1 , T p , T L ) ,    if     R est > 1 , where X est and R est respectively represent an expected bucket fill level and an estimated number of data units remaining to be transmitted in the message for the given connection. The most preferred embodiment relates to an ATM communications device wherein the messages are ATM adaption layer (AAL) service data units (SDUs) associated with a plurality of virtual circuit connections, and the data units are AAL segmentation and reassembly protocol data units (SAR PDUs). The transmission equipment preferably comprises segmentation means for segmenting each AAL SDU into one or more SAR PDUs, ATM layer means for multiplexing the SAR PDUs of various connections into a single stream having an output transmission rate, and physical transmission means, connected to the ATM layer means, for transporting the single stream of SAR PDUs over a physical interface. Another aspect of the invention relates to a method for transmitting multiple messages associated with multiple variable bit rate connections over a network, wherein each such connection is associated with a traffic contract which defines compliance thereto as conformance with a leaky bucket algorithm. The method comprises the steps of: (a) segmenting each message of each connection into one or more transportable data units; (b) multiplexing the data units from various connections into a single stream for transport over a transmission media to a destination; (c) scheduling the transfer of messages from the multiple connections to the multiplexing step; and (d) allocating a portion of the transmission media bandwidth to each connection substantially when its message is transferred to the multiplexing step, said portion of the transmission media bandwidth for a given connection being substantially equal to 1/T, T being computed as T ← max  ( T S + X - τ S N - 1 , T p , T L ) , if        N > 1 and T← max( T p , T L ), if N =1 where T S is a period corresponding to a constant sustained transmission rate, T p is a period corresponding to a peak transmission rate, τ S is a burst tolerance, N is the number of data units in the message, X is a fill level of the leaky bucket associated with the given connection, and T L corresponds to all unused bandwidth allocated to a service class to which the given connection belongs. Another aspect of the invention relates to a method for transmitting multiple messages associated with multiple variable bit rate connections over a network, wherein each such connection is associated with a traffic contract which defines compliance thereto as conformance with a leaky bucket algorithm. The method comprises the steps of: (a) segmenting each message of each connection into transportable data units; (b) multiplexing the data units from various connections into a single stream for transport over a transmission media to a destination; (c) scheduling the transfer of messages from the multiple connections to the multiplexing step; and (d) wherein the scheduling step includes one of: (i) scheduling the transfer of groups of data units from the segmentation step to the multiplexing step, wherein each such group of data units corresponds to a message, and (ii) scheduling the transfer of messages from the multiple connections to the segmentation step, wherein the segmentation step substantially immediately transfers the data units to the multiplexing step; and wherein the scheduling step is enabled to transfer the messages according to an order corresponding to the level of a leaky bucket associated with each corresponding connection in order to reduce the interleaving of data units from various connections over said single stream. BRIEF DESCRIPTION OF DRAWINGS For the purposes of description, but not of limitation, the foregoing and other aspects of the invention are explained in greater detail with reference to the accompanying drawings, wherein: FIG. 1 is an illustration of an ATM protocol stack; FIG. 2 is a protocol diagram illustrating the sub-layers of an ATM Adaption Layer (AAL) shown in FIG. 1; FIG. 3 is a data flow diagram illustrating the flow of PDUs through an ATM Layer shown in FIG. 1 using a round-robin scheduling discipline for servicing parallel instances or invocations of the AAL; FIG. 4A is a bandwidth occupancy chart illustrating the average latency in the transmission of SDUs from two virtual circuits using the round-robin scheduling discipline depicted in FIG. 3; FIG. 4B is a bandwidth occupancy chart exemplifying the average latency in the transmission of SDUs from two virtual circuits using apparatus according to one aspect of the invention; FIG. 5 is a protocol diagram illustrating the structure of the AAL in accordance with a first preferred embodiment of the invention which includes an SDU scheduling sublayer; FIGS. 6A and 6B are flowcharts illustrating a first preferred method, carried out by the scheduling sublayer, for servicing SDU data requests; FIG. 7 is a diagram illustrating various queues employed by the scheduling sublayer according to the first preferred embodiment; FIG. 8 is a table exemplifying average latency in transmitting SDUs using the first preferred method and using a prior art round-robin multiplexing scheme; FIG. 9 is a bandwidth occupancy chart corresponding to the example in FIG. 8 where SDUs are transmitted using the first preferred method; FIG. 10 is a bandwidth occupancy chart corresponding to the example in FIG. 8 where SDUs are transmitted using the prior art round-robin multiplexing scheme; FIG. 11 is a flow chart of a leaky bucket algorithm; FIG. 12 is a bandwidth occupancy chart exemplifying a transmission pattern for non-real-time VBR VCs in accordance with the first preferred method; FIG. 13 is a bandwidth occupancy chart exemplifying a transmission pattern for the same non-real-time VBR VCs shown in FIG. 12 but in accordance with a second preferred method for scheduling SDUs; FIGS. 14A and 14B are flowcharts illustrating the second preferred method, carried out by the scheduling sublayer, for servicing SDU data requests; FIG. 15 is a bandwidth occupancy chart illustrating a range of possible transmission times for PDUs of an SDU associated with a non-real-time VBR VC; FIG. 16 is a bandwidth occupancy chart exemplifying a transmission pattern for non-real-time VBR VCs in accordance with the second preferred method; and FIG. 17 is a bandwidth occupancy chart exemplifying a transmission pattern for the same non-real-time VBR VCs shown FIG. 16 but in accordance with the first preferred method. DETAILED DESCRIPTION OF PREFERRED EMBODIMENTS The detailed description is divided into two parts. The first part discloses a preferred embodiment for minimizing the average latency per unit of AAL SDU data in relation to ATM service classes such as the available bit rate (ABR) and unspecified bit rate (UBR) service categories which are characterized by a peak cell rate (PCR) traffic management parameter. The second part discloses a preferred embodiment for minimizing the average latency per unit of AAL SDU data in relation to ATM service classes such as non-real-time VBR which are characterized by peak cell rate (PCR), sustained cell rate (SCR), and maximum burst size (MBS) traffic management parameters. As will become apparent, the second embodiment builds upon the first embodiment. Service Classes Characterized by PCR FIG. 5 illustrates the structure of an AAL in accordance with the preferred embodiment. In addition to the known SAR layer 12 and CS layer 13 (which is not explicitly shown in FIG. 3 ), the preferred AAL includes a scheduling means, such as scheduling sublayer 24 (alternatively referred to as “SS”), located above the CS/SAR layers for scheduling or ordering the transfer of SDUs to the CS/SAR. (In alternative embodiments, scheduling sublayer 24 may be placed above the SAR but below the CS.) The scheduling sublayer 24 provides a plurality of SS SAPs 26 , one for each VC associated with PHY SAP 14 . The scheduling sublayer 24 accepts AAL SDUs 22 in data requests from higher layers of the ATM protocol stack (hereinafter “AAL user”). In the preferred embodiment, the AAL user uses the same data request to the scheduling sublayer at SS SAPs 26 as the AAL user would in the prior art for direct access to the AAL SAR SAP 15 . The scheduling sublayer 24 , in turn, preferably communicates with the SAR 12 using prior art AAL SAR SAP 15 . As will be described in greater detail below, the scheduling sublayer 24 affects only the egress data path (from AAL user to PHY) of SDUs 22 , but not the ingress data path (from PHY to AAL user). The scheduling sublayer 24 also communicates directly with the ATM Layer 11 through a cell rate specification signal 28 and an SDU transmission signal 30 . The cell rate signal 28 enables the scheduling sublayer 24 to specify a constant cell rate or period for any given VC to conventional ATM Layer equipment 11 ′ which, as known per se, provides this capability. The SDU transmission signal 30 enables the ATM Layer equipment 11 ′ to inform (such as by way of a software interrupt) the scheduling sublayer 24 when a last PDU of any given SDU 22 and VC has been transmitted. This capability is also provided by conventional ATM Layer equipment 11 ′, as known in the art per se, based upon examination of certain control bits in the PDU or when the equipment no longer has any cells to send for a particular connection. The scheduling sublayer 24 is provided with data regarding the traffic parameters, such as the PCR, and current transmit state of each VC; information which is typically known by the ATM Layer but not the CS/SAR. The scheduling sublayer 24 is also provided with data regarding SDU boundaries; information which is typically known by the CS/SAR but not the ATM Layer. This data is collectively used by the scheduling sublayer 24 to control the time at which SDUs 22 are passed to the CS/SAR, and to specify a constant cell rate for a given VC at the ATM Layer 11 . The scheduling sublayer 24 attempts to minimize the effects of the ATM Layer 11 which interleaves PDUs associated with various VCs in round-robin fashion for transport through PHY SAP 14 . This is accomplished by transferring the AAL SDUs of VCs to the CS/SAR layers on a substantially relative first-come, first-served basis (of VCs making data requests at SS SAPs 26 ) subject to the constraints that: (a) the PHY bandwidth consumed by the VCs are maximized within the limits of the bandwidth preallocated to their respective service classes; (b) the traffic parameters for each VC are satisfied; and (c) the SDUs, and the PDUs thereof, of any given VC be transmitted across PHY SAP 14 in sequential order so as to prevent misinsertion at the destination. This approach attempts to minimize the average transmission latency encountered by the AAL SDUs 22 . The potential results in accordance with this approach are simplistically illustrated in FIG. 4B for the situation, corresponding to FIG. 4A, wherein two VCs each request the transmission of one SDU having the same number of PDUs at about the same time, with each VC having a PCR equal to 100% of the available bandwidth. As shown in FIG. 4B, each VC is given a maximal amount of bandwidth, i.e., up to its PCR, but the average latency is reduced compared to FIG. 4 A. FIGS. 6A and 6B illustrate in flow chart form the preferred scheduling method carried out by the scheduling sublayer 24 . FIG. 7 illustrates various queues employed by the scheduling sublayer 24 . A TX list or queue 34 stores a list of VCs which are currently being transmitted by the CS/SAR and ATM Layer; that is, those VCs which have had an SDU (i.e., all of the PDUs thereof) passed to SAR 12 by the scheduling sublayer 24 . Queue 34 also stores the current cell rate for each VC listed therein. A TX-pending queue 36 stores a list of all VCs which have made data requests to the scheduling sublayer 24 , but the data requests and SDUs thereof have not yet been passed by the scheduling sublayer to AAL SAR SAP 15 . In addition, for each VC listed in queue 36 , there exists one VC buffer or queue 38 A, 38 B, 38 C, 38 D, . . . or 38 N for storing unprocessed SDUs of the corresponding VC. There are two triggering events in the preferred method. Event 40 (FIG. 6A) is the occurrence of a data request from the AAL user at any SS SAP 26 . Event 60 (FIG. 6B) is the communication of signal 30 by ATM Layer equipment 11 ′ to the scheduling sublayer 24 indicating that the equipment has finished transmitting an SDU (i.e., the last PDU thereof has been transmitted) for a given VC. Referring to FIGS. 5, 6 A and 7 , when the scheduling sublayer 24 receives at any SS SAP 26 a data request and the SDU thereof for a given VC, VC(i), the scheduling sublayer 24 checks at step 42 whether or not VC(i) is currently being serviced by it, i.e., if VC(i) is in the TX queue 34 or TX-pending queue 36 . If so, at step 46 , the scheduling sublayer 24 stores the SDU in a corresponding VC i queue 38 . These steps ensure that the AAL SDUs of a given VC are transferred to the CS/SAR in sequential order. If the condition at step 42 is not satisfied, then a check is made at step 44 whether or not there is any unused PHY bandwidth for the service class that VC(i) belongs to. If PHY bandwidth is available, then at steps 48 , 50 , 52 and 54 the scheduling sublayer 24 : (a) instructs the ATM Layer to set the current cell rate for VC(i) to the lower of: (i) the PCR of VC(i), and (ii) the remaining unused PHY bandwidth for the service class that VC(i) belongs to; (b) decreases the available PHY bandwidth for the service class that VC(i) belongs to by its current cell rate; (c) passes the SDU to the CS/SAR which, practically speaking, very quickly completes the segmentation of the SDU into PDUs and passes these to the ATM Layer for transmission; and (d) adds VC(i) to the tail end of the TX queue 34 . If, however, no PHY bandwidth is available at step 44 , at step 56 VC(i) is added to the tail end of the TX-pending queue 36 and the SDU of VC(i) is added to the tail end of the corresponding VC i queue 38 . Referring to FIGS. 5, 6 B and 7 , when the scheduling sublayer 24 receives an indication at event 60 that the ATM Layer 11 has finished sending all PDUs corresponding to an SDU of a given VC, VC(k), through the PHY SAP 14 , the SS 24 proceeds at step 62 to remove VC(k) from the TX queue 34 and update the PHY bandwidth availability for the service class that VC(k) belongs to. At step 64 , a check is made whether or not VC(k) has any AAL SDUs pending transmission, that is, if any SDUs are present in VC K queue 38 . If so, at step 66 VC(k) is placed on the TX -pending queue 36 . Steps 68 , 72 and 74 set up a loop for scanning the TX queue 34 , starting from its head. Step 70 determines whether the current cell rate of any VC, VC(x), in queue 34 can be increased. If so, at step 76 the scheduling sublayer 24 signals ATM Layer 11 to increase the current cell rate for VC(x) to the lower of: (i) the PCR of VC(x), and (ii) the remaining unused PHY bandwidth for the ATM service class that VC(x) belongs to. The available PHY bandwidth for that service class is also decreased accordingly. The TX queue 34 is scanned until the end of the queue or until the total PHY bandwidth available at PHY SAP 14 is exhausted. Steps 78 , 82 and 84 set up a loop for sequentially scanning the TX-pending queue 36 , starting from its head. If at step 80 PHY bandwidth remains available in the service class of a given VC, VC(x), at the head of queue 36 , then at step 86 the scheduling sublayer 24 removes VC(x) from TX-pending queue 36 . The scheduling sublayer 24 also removes the AAL SDU at the head of the corresponding VC X queue 38 , and processes the removed SDU as if a data request for VC(x) was received at the SS SAP 26 , as per steps 48 through 58 discussed above. The TX-pending queue 36 is scanned until the end of the queue or until the total PHY bandwidth available at PHY SAP 14 is exhausted. Referring to FIGS. 8, 9 and 10 , an example of how the scheduling sublayer 24 provides a lower average latency in transmitting SDUs than the prior art system shown in FIG. 3 . FIG. 8 is a table which specifies the VC identifier related to each data request, the time at which the request is made, the number of PDUs in the SDU, and the PCR for the VC. For simplicity, the VCs are assumed to belong to the same service class and the PCR is given as a percentage of the PHY bandwidth. The time has also been unitized so that one time unit is the length of time required for the PHY to send a single cell. FIG. 9 is a bandwidth occupancy chart illustrating the latency of each SDU that would likely result from the preferred scheduling method. This latency is the difference between the time a data request for the SDU was first made at SS SAP 26 and the time at which the last PDU for the SDU was sent by the ATM Layer to the PHY. For comparison, FIG. 10 illustrates the latencies which would result under a typical prior art round-robin ATM Layer multiplexing scheme. The average latencies per unit of SDU data,—that is, per PDU—are calculated and tabulated in the FIG. 8 table. The average latencies are calculated by taking the sums of the latency of each SDU multiplied by the number of PDUs in the SDU, and dividing by the total number of PDUs in all SDUs. Under the prior art system as shown in FIG. 10, all VCs begin transmission simultaneously, with each VC occupying 25% of the PHY bandwidth. At t=8, the SDU requested by VC 2 will have completed transmission. Accordingly, the prior art system equally reallocates the PHY bandwidth formerly occupied by VC 2 to VC 1 and VC 4 ; VC 3 is already transmitting at its PCR and cannot be allocated more bandwidth. The SDUs requested by VC 1 and VC 4 complete transmission at approximately time t=10.5, (in practice, one at time t=10 and one at time t=11). At t=16, the SDU requested by VC 3 completes transmission. The average latency is computed as (10.5*3+8*2+16*4+10.5*3)/(3+2+4+3)≈11.9. Under the preferred method, as shown in FIG. 9, at t=0, VC 1 is allocated its PCR, and the ATM Layer is instructed to send the PDUs of VC 1 at 50% of the PHY bandwidth. The SDU for VC 1 is passed to the SAR. Similarly, VC 2 is allocated its PCR, and the ATM Layer is instructed to send the PDUs of VC 2 at 33% of the PHY bandwidth, and the SDU for VC 2 is passed to the SAR. VC 3 is allocated the remaining PHY bandwidth, and the SDU for VC 3 is passed to the SAR. VC 1 , VC 2 and VC 3 are added to the TX queue 34 of VCs whose SDU transmission is in progress. Finally, since there is no remaining PHY bandwidth for VC 4 , the scheduling sublayer adds VC 4 to the TX-pending queue 36 of VCs which have data requests pending. This state is also illustrated in FIG. 7 A. At t=6, VC 1 and VC 2 both complete transmission of all PDUs for their SDUs. The scheduling sublayer 24 removes these VCs from the TX queue 34 of VCs whose transmission is in progress, and the PHY is marked as having 83% of its bandwidth previously allocated to those VCs as free. The TX-pending queue 36 of VCs whose transmission is in progress is then scanned, and VC 3 is found to be operating below its PCR, so the ATM Layer is instructed to increase the PHY bandwidth allocated to VC 3 from 17% to 25%. Since there is still free PHY bandwidth, the TX-pending queue 36 of VCs with data requests pending is scanned. VC 4 is found, and the ATM Layer is instructed to allocate 75% of the PHY bandwidth to VC 4 —75% being both VC 4 's PCR and the previously remaining free PHY bandwidth - and the SDU for VC 4 is passed to the SAR. VC 4 is then added to the TX queue 34 of VCs whose transmission is in progress. This state is also illustrated in FIG. 7 B. At t=10, VC 4 completes transmission of PDUs for its SDU. VC 4 is therefore removed from the TX queue 34 of VCs whose transmission is in progress. Queue 34 is then scanned. VC 3 is found, but it is already transmitting at its PCR, so its transmission is not modified. There are no other VCs with data requests pending, so the PHY becomes 75% idle. At t=18, VC 3 completes transmission. The average latency under the scheduling sublayer 24 is computed as (6*3+6*2+18*4+10*3)/(3+2+4+3)=11. The preferred embodiment provides a number of advantages. First, it schedules SDUs to minimize the average amount of latency per unit of AAL SDU data on a given PHY. This results in effectively a lower average delay in data transmission through an ATM network. Second, the preferred method features a low storage and computational load, and can therefore be easily applied in real time systems. Third, the preferred embodiment strives to ensure that all available PHY bandwidth is consumed. Fourth, the preferred embodiment is not dependent on elaborate ATM Layer equipment. The only requirement of the ATM Layer equipment is that it be configurable so that the transmission rate of each connection or VC can be dynamically adjusted. However, while the preferred embodiment employs ATM Layer equipment capable of generating a direct indication, e.g. software interrupt, of when an AAL SDU has finished being transmitted through the ATM Layer, alternative embodiments need not provide this capability. Instead, the scheduling sublayer may use a separate timer to provide this indication. This is because the scheduling sublayer 24 is provided with the cell rate of each VC and the number of AAL PDUs to be sent on each connection. Therefore, the scheduling sublayer can compute the approximate time (i.e., ignoring cell delay variation introduced by ATM Layer multiplexing and PHY functions) at which the ATM Layer will finish sending the ATM PDUs to the PHY SAP. It should also be realized from the foregoing that the scheduling sublayer 24 implements a first-come, first-served scheduling discipline with respect to data requests at the point of contention. For example, referring to FIGS. 8 and 9, VCs 1 , 2 , 3 and 4 submit SDU data requests in that order, so the scheduling sublayer 24 services the connections in that order. In the particular example shown, the SDU of VC 4 cannot be immediately transferred to the CS/SAR because no bandwidth is available due to the servicing of the other VCs and so the SDU of VC 4 is forwarded to the CS/SAR at time t=6. Nevertheless, through the use of queues 34 , 36 and 38 , the scheduling sublayer 24 applies a round-robin service discipline with respect to the connections per se. For instance, referring to FIGS. 8 and 9, suppose that at time t=0 + , VC 1 presents a second SDU (i.e., two SDUs) to the scheduling sublayer 24 before VC 4 presents its SDU thereto. Due to step 42 of the preferred method, the scheduling sublayer 24 places the second SDU of VC 1 in VC 1 queue 38 A. The SDU of VC 4 is placed in VC 4 queue 38 D, but VC 4 is itself placed in the TX-pending queue 36 . At time t=6 when the first SDU of VC 1 has finished transmitting, step 64 of the preferred method notes that VC 1 has an outstanding data request and thus places VC 1 on the TX-pending queue 36 , behind VC 4 . Thus, at time t=6 when bandwidth becomes available, the (only) SDU of VC 4 is submitted to the CS/SAR by step 86 ahead of the second SDU of VC 1 , even though the latter physically arrived at the scheduling sublayer 24 prior to the former. In this manner, the preferred method provides a means for ensuring that no VC is starved. In alternative embodiments, the scheduling sublayer could implement a strict first-come first-served order per SDU, e.g., by transferring each SDU to the CS/SAR in the order of arrival as follows: (i) immediately, provided that bandwidth is available, or in due course, when bandwidth becomes available to the VC of that SDU or (ii) immediately, if the VC of that SDU is currently transmitting an SDU, thereby essentially changing the queuing point for such VCs to the ATM Layer. Such an embodiment may, however, risk starving a connection. The AAL according to preferred embodiment as shown in FIG. 5 places the scheduling sublayer 24 above the CS/SAR layers 12 and 13 . It should be appreciated that in alternative embodiments the scheduling sublayer 24 may be placed between the SAR layer 12 and the ATM layer 11 . In such an embodiment, the scheduling sublayer transfers groups of AAL SAR PDUs to the ATM layer 11 , wherein each group of AAL SAR PDUs corresponds to an AAL SDU. This embodiment yields an AAL having the same functionality as that of the preferred embodiment because the payload of the group of AAL SAR PDUs is in effect the AAL SDU. Service Classes Characterized by PCR and SCR The embodiment described above operates well when the ATM virtual connections have traffic parameters such that they are always permitted to transmit at rates up to some peak cell rate. However, further optimization of the scheduling sublayer 24 can be obtained for the non-real-time variable bit rate (nrt-VBR) service category, where the rate at which a VC can transmit at a particular instant may vary between a specified sustainable cell rate and a peak cell rate. ATM Forum document no. AF-TM-0056.000 and ITU specification no. I.371 describe the VBR traffic contract in terms of peak cell rate (PCR), sustainable cell rate (SCR) and maximum burst size (MBS). Compliance to the VBR traffic contract is defined in terms of the generic cell rate algorithm (GCRA), also commonly known as the “leaky bucket” algorithm, which is shown in FIG. 11. A VC is said to be compliant if the cell stream carried by that VC conforms to both GCRA (T S , τ S ) and GCRA (T p , 0), where T S is the cell period at the SCR, ( T S = 1 SCR ) , T p is the cell period at the PCR, ( T p = 1 PCR ) , and τ S is the burst tolerance, (τ S =(MBS−1)*(T S −T p )). The essence of the non-real-time VBR traffic contract is that it explicitly allows a VC to transmit in a “bursty” fashion. That is, a nrt-VBR VC can transmit indefinitely at its SCR, and be in conformance with its traffic contract. Alternatively, the VC can cyclically transmit at rates above its SCR, up to its PCR, for a well-defined period, followed by another period during which it must transmit at rates below its SCR. This bursty transmission pattern is also in conformance with the nrt-VBR traffic contract. FIG. 12 is a bandwidth occupancy chart illustrating a transmission pattern that arises under the scheduling sublayer 24 of the first preferred embodiment for two nrt-VBR VCs. Each of these VCs has a set of AAL SDUs to send, with traffic parameters specifying an SCR of 50% of the PHY bandwidth and a PCR of 100% of the PHY bandwidth. It will be seen that each VC is continually sending at its SCR, and there is no burstiness in transmission. This transmission pattern is in conformance with the example nrt-VBR traffic contract, and the average latencies are noted on the chart. The second preferred embodiment improves upon the first preferred embodiment by modifying the scheduling sublayer 24 to schedule the transmission of an SDU such that the burstiness permitted by the nrt-VBR traffic contract is advantageously used to further reduce the average latency in the transmission of an SDU. This is done by scheduling the transmission of SDUs from various VCs such that the VCs may alternate between transmitting at rates approaching their PCRs and pausing their transmission, and by having the bursts in the transmission of each VC contain exactly one SDU. The transmission pattern under this approach is shown in FIG. 13 for the same example traffic contract described with reference to FIG. 12 . It will be seen that while one VC has paused its transmission, the other VC is sending at its PCR. This transmission pattern is highly bursty, but also conforms to the example nrt-VBR traffic contract. As in the first embodiment, the second embodiment employs ATM Layer equipment 11 ′ which preferably has the following characteristics: Each VC whose cells are multiplexed by the ATM Layer may be assigned its own constant cell period. This cell period may be modified at any time by some other entity in the system, such as the scheduling sublayer 24 , but it is not otherwise modified by the ATM layer. The ATM Layer performs cell multiplexing such that, for each opportunity it has to send a cell to the PHY, it may select a cell from any of the VCs for which it has candidate cells available for transmission, provided only that the interval between the current instant and the time at which the previous cell was transmitted for a particular VC is not less than the cell period specified for that VC. In the event that the cell period specified for a VC is modified by some entity such as the scheduling sublayer, the interval between the time at which the last cell was transmitted at the old cell period and the first cell is transmitted at the new cell period is not less than the lesser of the old and new cell periods. FIGS. 14A and 14B show the preferred method carried out by the scheduling sublayer 24 in accordance with the second preferred embodiment. Those steps which have been added or modified in relation to FIGS. 6A and 6B of the first embodiment are shown in stippled outlines. In the second preferred method, the scheduling sublayer 24 stores two parameters, t and X, for each VC. X corresponds to the “fill level” of the GCRA(T S , τ S ) “leaky bucket”, measured at time t. Initially, i.e., when a given VC is established, both X and t are zero. X and t are preferably recomputed by the scheduling sublayer 24 for each VC at two instants: (i) at a step 47 a (FIG. 14 A), before the scheduling sublayer 24 passes (at step 52 ) a given SDU 22 to the CS/SAR 12 , and (ii) at a step 61 (FIG. 14 B), after the scheduling sublayer 24 receives (at event 60 ) an indication that the transmission of the given SDU 22 is complete. At step 47 a , before the scheduling sublayer 24 requests the SAR 12 to transmit a given SDU, X and t are calculated as follows: X← max( X−t now +t , 0) t←t now   {1} where t now represents the time at which the calculation is performed. This calculation represents the “draining” of the leaky bucket in the period from t to t now , during which there was no transmission of PDUs for the VC. At step 61 , after the scheduling sublayer 24 receives an indication that the transmission of the given SDU is complete, X and t are calculated as follows:   X← max( X−t now +t+N*T S , 0) t←t now   {2} where N is the number of PDUs (that is, the number of ATM cells) into which the given SDU was segmented by the SAR 12 . This calculation represents the “filling” of the leaky bucket in the period from t to t now , during which N PDUs were transmitted on the VC. When the scheduling sublayer 24 has a candidate SDU which could begin transmission immediately, it checks whether the first cell of that SDU would be conformant with GCRA(T S , τ S ) if the SDU were transmitted immediately, i.e., if X−t now +t≦τ S . Such a check occurs at step 47 b (FIG. 14 A), which corresponds to the condition where an AAL user sends a data request to the scheduling sublayer 24 , PHY bandwidth is currently available, and queues 34 and 36 of pending SDUs for that VC are empty. This check also preferably occurs at step 81 a (FIG. 14B) which corresponds to the situation where another SDU has completed transmission and the scheduling sublayer is scanning the TX-pending queue 36 of VCs with pending SDUs. If the check, X−t now +t≦τ S , is satisfied, the scheduling sublayer 24 can pass the candidate SDU to the CS/SAR 12 for immediate transmission. If the check is not satisfied, the scheduling sublayer 24 cannot immediately begin transmitting the candidate SDU without violating GCRA(T S , τ S ). Instead, in the case of a newly submitted SDU, at step 56 ′ the SDU is placed at the end of the queue 38 of pending SDUs for that VC, and that VC is placed on the TX-pending queue 36 of VCs which have pending SDUs. In the case where the scheduling sublayer 24 is scanning the TX-pending queue 36 of VCs with pending SDUs, queue 36 is left untouched. In either case, at steps 47 c (FIG. 14A) and 81 b (FIG. 14 B), the scheduling sublayer 24 starts a timer set to a period X−τ S −t now +t. When this timer expires at event 79 (FIG. 14 B), the scheduling sublayer will again scan the TX-pending queue 36 of VCs with pending data requests. At that time, the former candidate SDU should be able to be transmitted in conformance with GCRA(T S , τ S ). Before the scheduling sublayer passes an SDU to the CS/SAR for transmission (step 52 ), at step 48 ′ the scheduling sublayer computes a cell period, T, (T being the inverse of the cell rate) which should be applied by the ATM Layer to the VC over the course of the transmission of the SDU. When N>1, T is calculated according to: T ← max  ( T S + X - τ S N - 1 , T p , T L )     if     N > 1 {3a} where T L represents the period at which the remaining bandwidth of the PHY would be completely consumed by the transmission of this SDU. The first term of the above max( ) function represents the cell period which would cause the “bucket” of the VC to be completely filled at the end of the SDU's transmission. This implies that the SDU will be transmitted in as fast a burst as possible while still being in conformance with GCRA(T S , τ S ). (Note that X=τ S implies that the bucket is full when the SDU's transmission begins, in which case the first term of the max( ) function simplifies to T S —when the bucket is full, the cell period cannot be less than the sustainable cell period.) The second term prevents the cell period from being less than the peak cell period. The third term prevents the VC from attempting to consume more PHY bandwidth than exists. When N=1, the first term in equation {3a} is not considered. Thus,   T← max( T p , T L )  {3b} After an SDU's transmission is completed (event 60 ), the scheduling sublayer 24 attempts at step 70 ′ to increase the cell rate, i.e., decrease the cell period, of any VCs with SDUs currently in transmission. To accommodate this, the scheduling sublayer 24 stores parameters X est ,R est and t est for each VC. These represent the estimated bucket fill level, the estimated number of cells remaining to transmit in the SDU, and the time of estimation, respectively. They are initialized at step 52 , when the scheduling sublayer 24 passes an SDU to the CS/SAR for transmission, according to: X est ←X, R est ←N, t est ←t   {4} Thereafter, if the VC's cell period is adjusted while the SDU is in transmission, the new cell period is computed at step 70 ′ according to: X est ← X est + ( t now - t est )     T S - T T {5a} R est ← R est - t now - t est T {5b}   t est ←t now   {5c} T ← max  ( T S + X est - τ S R est - 1 , T p , T L ) ,    if     R est > 1 {5d} If R est ≦1, the cell period of the VC is not modified. Again, the first term in the max( ) of equation {5d} results in a cell period T such that the bucket will be full when the SDU's transmission is complete. The “est” subscript in the above expressions refer to the fact that X est and R est are estimates of the bucket fill level and number of cells remaining to be transmitted in the SDU, based on the elapsed time since the last estimate (or since the SDU's transmission was requested to the SAR) and the cell period; they are not based on direct indications from the ATM layer or the SAR. In the event that the ATM Layer can provide a direct indication of the number of cells remaining to be transmitted, then the scheduling sublayer need not track X est , separately from X. Alternatively, in the event that the ATM Layer does not provide direct indications at any time, including when the last cell of an AAL SDU has been transmitted, there is again no need to distinguish between X est and X. X est should be distinguished from X only if the calculations during the transmission of an SDU are time based estimates and the calculations at the end of the transmission of an SDU are based on a direct indication from the ATM Layer. The advantage of this distinction is that any errors which may accumulate in the estimations are discarded when the direct indication is received at the ATM Layer. The foregoing additions and modifications to the first preferred embodiment have the objective of determining the maximum possible cell rate which can be used for a particular VC at a particular instant in time while maintaining conformance to the VBR traffic contract. In addition to the foregoing additions and modifications, at steps 56 ′ (FIG. 14A) and 66 ′ (FIG. 14B) the scheduling sublayer 24 sorts the TX-pending queue 36 of VCs which have SDUs pending transmission according to increasing values of parameter t E : ti t E ←X+t   {6} t E is simply the time at which the VC's bucket will be empty. This means that when one SDU's transmission has completed and another SDU's transmission may begin, the SDU selected to begin transmission will be the one on the VC whose bucket is emptiest. The VC with the emptiest bucket is that which can transmit at a cell rate closest to its PCR, and is therefore the VC which can achieve the most bursty transmission. Steps 56 ′ and 66 ′ attempt to maximize burstiness by taking advantage of the linearity of the GCRA. As shown in FIG. 15, there is a range of times when transmission of an SDU can begin which result in the same time for the completion of transmission of the last cell for that SDU, and hence the same latency therefor. Within that range, if the start of transmission of an AAL SDU is delayed such that the VC is not transmitting during that delay period, the VC's bucket will be emptying. This allows the SDU to eventually be transmitted at a cell rate approaching the PCR, which recoups the delay time. Thus, the increase in burstiness has no impact on that latency of the SDU. However, in the idle time before that transmission of that SDU begins, it is possible that another SDU could be completely transmitted on another VC, thereby lowering the average SDU latency of the system. An example is now presented of the advantages obtained by sorting the TX-pending queue 36 in steps 56 ′ and 66 ′ in order to maximize burstiness and hence reduce latency. Consider a system having three nrt-VBR VCs with each VC having the same traffic parameters: a PCR of 100% of the PHY bandwidth, an SCR of 50% of the PHY bandwidth, and an MBS of 4. This yields T p =1, T S =2 and t S =3 for each VC. (As before, a unit of time is deemed to be the time required to transmit one SAR PDU or cell at the PHY.) Suppose at time 0 , the scheduling sublayer 24 receives a data request on VC 1 for an SDU which will be segmented into 4 SAR PDUs, and immediately thereafter receives a data request on VC 2 for an SDU which will fit into a single SAR PDU. Just after time 4 , the scheduling sublayer receives another data request on VC 1 for an SDU which will be segmented into 3 SAR PDUs. Also just after time 4 , the scheduling sublayer receives a data request on VC 3 for an SDU which will be segmented into 3 SAR PDUs. Assume that all VCs are in their initial state, with empty buckets, at time 0 . Referring additionally to the bandwidth occupancy chart of FIG. 16, the scheduling sublayer proceeds as follows: t now =0: First, the scheduling sublayer verifies that VC 1 is GCRA(T S , τ S ) compliant, i.e., that X−t now +t≦τ S for VC 1 . Since this condition is met, VC 1 can begin transmission of an AAL SDU immediately. The scheduling sublayer then computes (X, t)=(0, 0) for VC 1 . The period at which PDUs should be transmitted for this SDU is, from equation {3a}, given by: T = max  ( T S + X - τ S N - 1 , T p , T L ) = max  ( 2 + 0 - 3 4 - 1 , 1 , 1 ) = 1 In other words, the scheduling sublayer will tell the ATM Layer to set the cell rate of VC 1 to 100% of the PHY bandwidth, and will request the transmission of the SDU for VC 1 to the SAR. Also, since all the PHY bandwidth has been allocated at this time, the SDU requested for VC 2 is added to the end of the VC 2 queue 38 b of pending SDUs for VC 2 , and VC 2 is placed on the TX-pending queue 36 for VCs having pending data requests, with t E =0. t now =4: The transmission of the SDU for VC 1 is completed. The parameters (X, t)=(0−4+0+2*4, 4)=(4, 4) are stored against VC 1 . The transmission latency for this SDU is 4. There are no other VCs with SDUs in transmission. The TX-pending queue 36 of VCs having pending data requests is scanned, finding VC 2 . For VC 2 , the condition X−t now +t≦τ S is met, so transmission may begin. The parameters (X, t)=(0, 4) are computed for VC 2 . Since N=1 for the SDU requested on VC 2 , the transmission of the SDU occurs at the peak cell period of 1. The scheduling sublayer 24 thus signals the ATM Layer 11 to set the cell rate of VC 2 to 100% of the PHY bandwidth, and will request the transmission of the SDU for VC 2 to the SAR 12 . t now =4 + : Next, the scheduling sublayer receives the data request for the second SDU on VC 1 . Since there is no unused PHY bandwidth remaining, the SDU requested for VC 1 is added to the end of VC, queue 38 a of pending SDUs for VC 1 , and VC 1 is placed on the TX-pending queue 36 for VCs with pending SDU data requests, with t E =X+t=8. Finally, the scheduling sublayer receives the data request for the SDU of VC 3 . Since there is no unused PHY bandwidth remaining, the SDU for VC 3 is added to the end of the VC 3 queue 38 c of pending SDUs for VC 3 , and VC 3 is placed on the TX-pending queue 36 of VCs with pending SDUs. Since the t E value for VC 3 is 0, which is lower than the t E value for VC 1 (t E =8 for VC 1 ), VC 3 is inserted before VC 1 on the TX-pending queue 36 for VCs having SDU data requests pending, in accordance with step 66 ′. t now =5: The transmission of the SDU for VC 2 is completed. The parameters (X, t)=(0−5+4+2*1, 5)=(1, 5) are stored against VC 2 . The transmission latency for this SDU is 2. There are no other VCs with SDUs in transmission. The TX-pending queue 36 of VCs having pending SDUs is scanned, finding VC 3 first. For VC 3 , the condition X−t now +t≦τ S met, so transmission may begin. The parameters (X, t)=(0, 5) are computed for VC 3 . The period at which PDUs should be transmitted for the SDU of VC 3 is given by: T = max  ( T S + X - τ S N - 1 , T p , T L ) = max  ( 2 + 0 - 3 4 - 1 , 1 , 1 ) = 1 The scheduling sublayer thus signals the ATM layer to set the cell rate of VC 3 to 100% of the PHY bandwidth, and will request the transmission of the SDU for VC 3 to the SAR. t now =8: The transmission of the SDU for VC 3 is completed. The parameters (X, t)=(0−8+5+2*3, 8)=(3, 8) are stored against VC 3 . The transmission latency for this SDU is 4. There are no other VCs with SDUs in transmission. The TX-pending queue 36 of VCs having pending SDU data requests is scanned, finding VC 1 . For VC 1 , the condition X−t now +t≦τ S is met, so transmission may begin. The parameters (X, t)=(4−8+4, 8)=(0, 8) are computed for VC 1 . The period at which PDUs should be transmitted for the SDU of VC 3 is given by: T = max  ( T S + X - τ S N - 1 , T p , T L ) = max  ( 2 + 0 - 3 4 - 1 , 1 , 1 ) = 1 The scheduling sublayer thus signals the ATM Layer to set the cell rate of VC 1 to 100% of the PHY bandwidth, and will request the transmission of the second SDU for VC 1 to the SAR. t now =11: The transmission of the second SDU for VC 1 is completed. The parameters (X, t)=(0−11+8+2*3, 11)=(3, 11) are stored against VC 3 . The transmission latency for this SDU is 7. The average transmission latency per unit of SDU data is then (4*4+2*1+4*3+7*3)/(4+1+3+3)≈4.6. Referring now to the bandwidth occupancy chart of FIG. 17, a contrasting example is now presented assuming the sorting steps 56 ′ and 66 ′ of the second preferred embodiment are not applied so that the TX-pending queue 36 of VCs with SDUs pending is maintained in simple first-in-first-out (FIFO) order as in the first preferred embodiment. In this case, after t now =4, the order of the TX-pending queue 36 of VCs having pending SDUs would have VC 1 before VC 3 . Accordingly, the following modifications would occur: t now =5: The transmission of the SDU for VC 2 is completed. The parameters (X, t)=(0−5+4+2*1, 5)=(1, 5) are stored against VC 2 . The transmission latency for this SDU is 2. There are no other VCs with SDUs in transmission. The TX-pending queue 36 of VCs with pending SDUs is scanned, finding VC 1 first. For VC 1 , the condition X−t now +t≦τ S is met, so transmission may begin. The parameters (X, t)=(4−5+4, 5)=(3, 5) are computed for VC 1 . The period at which PDUs should be transmitted for the second SDU of VC 1 is given by: T = max  ( T S + X - τ S N - 1 , T P , T L ) = max  ( 2 + 3 - 3 3 - 1 , 1 , 1 ) = 2 The scheduling sublayer thus instructs the ATM Layer to set the cell rate of VC 1 to 50% of the PHY bandwidth, and will request the transmission of the second SDU for VC 1 to the SAR. Fifty per cent (50%) of the PHY bandwidth remains unused, so the scanning of the TX-pending queue 36 of VCs with SDUs pending continues. VC 3 is found. For VC 3 , the condition X−t now +t≦τ S is met, so transmission may begin. The parameters (X, t)=(0, 5) are computed for VC 3 . The period at which PDUs should be transmitted for VC 3 's SDU is given by: T = max  ( T S + X - τ S N - 1 , T P , T L ) = max  ( 2 + 0 - 3 4 - 1 , 1 , 2 ) = 2 This time, the limitation of the bandwidth allocated to VC 3 is not imposed by GCRA conformance, but by the limited amount of PHY bandwidth remaining. The scheduling sublayer thus instructs the ATM Layer to set the cell rate of VC 3 to 50% of the PHY bandwidth, and will request the transmission of the SDU for VC 3 to the SAR. t now =11: The transmission of the second SDU for VC 1 and the SDU for VC 3 are both completed. The parameters (X, t)=(3−11+5+2*3, 11)=(3, 11) are stored against VC 1 ; the parameters (X, t)=(0−11+5+2*3, 11)=(0, 11) are stored against VC 3 . The transmission latency for both SDUs is 7. The average transmission latency per unit of SDU data is then (4*4+2*1+7*3+7*3)/(4+1+3+3)≈5.5—worse than is achieved when the queue of VCs with SDUs pending is sorted by increasing t E . The discussion now turns to mathematically demonstrating that the second preferred embodiment produces a cell stream conformant to the nrt-VBR traffic contract. The analysis refers to the algorithms and discussion presented in ITU specification I.371 and especially in ATM Forum document AF-TM-0056.000, Appendix C, and uses a modified version of the notation in the “continuous-state leaky bucket algorithm” presented as an equivalent of the GCRA in AF-TM-0056.000 and I.371. Let t K be the time at which the K th conforming cell was transmitted on an ATM virtual connection, let X′ K be the trial value of the “leaky bucket counter” before the K th conforming cell was transmitted, and let X K be the value of the “leaky bucket counter” of the GCRA(T S , t S ) after the K th conforming cell was transmitted. If a conforming cell is transmitted at t K , then the relationship between X′ K and X K is given by: X K =X′ K +T S   {7} Suppose that at time t′ K+1 , the SAR requests of the ATM layer that N ATM SDUs, that is, cells, be transmitted. The minimum values of t K+1 and t K+N —that is, the earliest time at which the first and last of the N cells can be transmitted—can be calculated such that compliance to the traffic contract is guaranteed. Further, it can be shown that if N>1, the N cells can be transmitted at some rate with constant period T. The linearity of the GCRA guarantees that if the first and last of N cells are conforming, and all cells between the first and last are equally spaced, then all cells between the first and last will also be conforming. Note that transmitting cells at a fixed rate with period T is but one of many possible transmission scenarios that would be compliant to a given traffic contract. First, the minimum value of t K+1 is calculated: t K+1 ≧max( t′ K+1 ,t K +T p ,t K +X K −τ S )  {8} The second term of the max( ) function in {8} assures compliance with GCRA(T p , 0); the last term assures compliance with GCRA(T S , τ S ). The first term is pro forma, and indicates that the ATM layer cannot send the first cell before the ATM layer user has requested it. Before the first cell is transmitted at t K+1 , the value of X′ K+1 is given by: X′ K+1 =max( X K −t K+1 +t K ,0)  {9} The first term of the max( ) function in {9} corresponds to the case where there has not been sufficient time between t K and t K+1 to “empty the bucket”; the second (zero) term of the max( ) function indicates that there has been sufficient time between t K and t K+1 to “empty the bucket”. Next, if N>1, assume there will be N−1 iterations through the GCRA after t K+1 up to and including t K+N , where each iteration is separated by equal time T. After these N−1 iterations, to conform to GCRA(T S , τ S ), T must satisfy:  ( N− 1)( T S −T )+X′ K+1 ≦τ S   {10} To conform to GCRA(T p , 0), T must satisfy T=T p . Therefore, combining these constraints on T, obtain for the minimum value of T: T ≥ max  ( T S + X K + 1 ′ - τ S N - 1 , T p ) { 11 } Similarly, after N−1 iterations through the GCRA, calculate X K+N , according to: X K+N =NT S −( N− 1) T+X′ K+1   {12} Note that {12} applies for all N, including N=1. For N−1, however, the need to calculate T is obviated by the multiplication of T by N−1. When this algorithm is first run at time t′ 1 for a particular VC, the values t 0 and X 0 are initialized to t 0 =t′ 1 and X 0 =0. It can be shown that relaxing the condition that the N cells be transmitted with equal period T does not change the minimum time t K+N , but the loss of constant T does cause the loss of the guarantee that cells between the first and Nth are conformant. Therefore, there is no benefit but there is some cost in spacing the N cells equally. Referring back to the operation of the second preferred embodiment, the scheduling sublayer's per-VC parameter X immediately before the scheduling sublayer requests transmission of an AAL SDU of the SAR, corresponds to X′ K+1 in this analysis; the parameter t corresponds to t K+1 in the analysis. Once the ATM layer has completed transmission of all SAR PDUs and has indicated this back to the scheduling sublayer, the scheduling sublayer's per-VC parameter X corresponds to X K+N in this analysis, the parameter t corresponds to t K+N in the analysis. Thus, {9} gives {1}, {12} gives {2}, {11} gives {3a}, and {8} gives the conditions stated in the second desired characteristic of the ATM layer equipment. The minimum time at which the last conformant SAR PDU of an AAL SDU can be transmitted by the ATM layer is then t K+N , where: t K + N =    t K + 1 + ( N - 1 )  T =    t K + 1 + ( N - 1 ) * max  ( T S + X K + 1 ′ - τ S N - 1 , T P ) =    t K + 1 + ( N - 1 ) * max  ( T S + max     ( X K - t K + 1 + t K , 0 ) - τ S N - 1 , T P ) { 13 } The expansions in {13} follow from {11} and {9}. Now, suppose the following inequalities are satisfied: X K −t K+1 +t K ≧0  {14a} T S + X K + 1 ′ - τ S N - 1 ≥ T P { 14  b } Then, {13} can be simplified to: t K + N =    t K + 1 + ( N - 1 )  ( T S + X K - t K + 1 + t K - τ S N - 1 ) =    ( N - 1 )     T S + X K + t K - τ S { 15 } The result of {15} means that t K+N is independent of t K+1 provided conditions {14a} and {14b} are met. The inequalities {8}, {14a} and {14b} can be combined with the definition of t S and rearranged as: X K +t k τ S ≦t K+1 ≦X K +t k +min(( N−MBS )( T S −T p ),0)  {16} This gives the range over which the time of transmission of the first SAR PDU for an AAL SDU may vary, without affecting the time of transmission of the last SAR PDU for the AAL SDU and hence the latency of the AAL SDU. In other words, this is the time range illustrated in FIG. 15 . Those skilled in the art will appreciate that numerous modifications and variations may be made to the preferred embodiments without departing from the spirit and scope of the invention.

The ATM apparatus for transmitting messages associated with a plurality of variable bit rate connections comprises ATM layer transmission equipment for receiving multiple messages, segmenting each received message into a plurality of smaller data units, such as 48 byte ATM adaption layer segmentation and reassembly protocol data units, and multiplexing such data units into a single stream for transport over a physical interface. A scheduler receives messages from each of the variable bit rate connections and transfers the messages to the transmission equipment in an order corresponding to the level of a leaky bucket associated with each connection. The scheduler also dynamically sets a transmission rate, 1/T, for each connection at the time its message is transferred to the transmission equipment, where T ← max  ( T S + X - τ S N - 1 , T p , T L ) , if        N > 1 and T← max( T p , T L ), if N =1 where T S is a period corresponding to a constant sustained transmission rate, T p is a period corresponding to a peak transmission rate, τ S is a burst tolerance, N is the number of data units in the message, X is a fill level of the leaky bucket associated with the given connection, and T L corresponds to all unused bandwidth allocated to a service class to which the given connection belongs. These actions increase the burstiness of the connections and minimize interleaving of data units from the various connections in order to reduce the average latency (per unit of message) in transmitting the messages across a network.

BACKGROUND The DNA genome in the particle of a DNA virus can be double-stranded, single-stranded, or partially double-stranded DNA. Papovaviridae (e.g., papillomaviruses), herpesviridae (e.g., herpes simplex viruses), and adenoviridae (e.g., adenoviruses) contain double-stranded DNA genomes. Some viruses from parvoviridae (e.g., parvoviruses) contain single-stranded DNA as their genomes. Some viruses from hepadnaviridae (e.g., hepatitis B virus) contain partially double-stranded DNA as their genomes, and replicate their genomes through RNA intermediates. In contrast, retroviruses are a family of RNA viruses that replicate through a DNA intermediate. Examples of retroviruses include Moloney murine sarcoma viruses, human T-cell lymphotrophic viruses, human immunodeficiency viruses, and human foamy viruses. Viruses described above cause a variety of diseases, such as the flu, common cold, herpes, measles, small pox, and encephalitis. Vaccination only offers protection for uninfected individuals for a few viral diseases. Thus, there is a need for identifying therapeutic agents useful for preventing or treating viral infection. SUMMARY The present invention is based, in part, on the discovery of nucleotide analogs that possess anti-viral activity. In one aspect, this invention relates purine compounds of formula (I): R 1 is NH 2 or OH; R 2 is H or NH 2 ; R 3 is H or alkyl; each of m and n, independently, is 1, 2, 3, or 4; X is O, S, or NH; and Y is H, halogen, OR a , P(O)(OR a ) 2 , or P(O)(OR a )(OR b ), in which R a is H, alkyl, aryl, heteroaryl, cyclyl, heterocyclyl, and R b is (referred to as “sugar-A(R c )—OP(O)(OR d )(OR e )” hereinafter), wherein A is adenine, guanine, cytosine, uracil, or thymine; R c is H or OH; R d is H or alkyl; R e is H, alkyl, or 5-ethylidene-(3,4-dialkoxyl)-furan-2-one; provided that if R 1 is NH 2 , R 2 is H; and if R 1 is OH, R 2 is NH 2 . Note that the left atom shown in any of the substituted groups set forth above is closest to the purine ring. Referring to formula (I), a subset of the purine compounds are those in which each of R 1 is NH 2 and R 2 is H. In these compounds, R 3 can be H; X can be O; and each of m and n, independently, is 1 or 2. In some embodiments, m is 1, X is O, n is 2, and Y is OR a , in which R a can be H. In other embodiments, m is 2, X is O, n is 1, and Y is P(O)(OR a ) 2 , in which R a can be H. In still other embodiments, m is 2, X is O, n is 1, and Y is P(O)(OR a )(OR b ), in which R a can be H, and R b can be sugar-A(R c )—OP(O)(OR d )(OR e ). Another subset of the purine compounds of formula (I) are those in which R 1 is OH and R 2 is NH 2 . In these compounds, R 3 can be H; X can be O; and each of m and n, independently, is 1 or 2. In some embodiments, m is 1, X is O, n is 2, and Y is OR a , in which R a can be H. In other embodiments, m is 2, X is O, n is 1, and Y is P(O)(OR a ) 2 , in which R a can be H. In still other embodiments, m is 2, X is O, n is 1, and Y is P(O)(OR a )(OR b ), in which R a can be H, and R b can be sugar-A(R c )-OP(O)(OR d )(OR e ). In another aspect, this invention encompasses purine compounds of formula (II): R 1 is NH 2 or OH; R 2 is H or NH 2 ; R 3 is H or alkyl; each of m and n, independently, is 1, 2, 3, or 4; X is O, S, or NH; and Y is P(O)(OR a )(OR b ), in which R a is H, alkyl, aryl, heteroaryl, cyclyl, heterocyclyl; and R b is 5-ethylidene-(3,4-dialkoxy)-furan-2-one or sugar-A(R c )—OP(O)(OR d )(OR e ); wherein A is adenine, guanine, cytosine, uracil, or thymine; R c is H or OH; R d is H or alkyl; R e is H, alkyl, or 5-ethylidene-(3,4-dialkoxyl)-faran-2-one; provided that if R 1 is NH 2 , R 2 is H; and if R 1 is OH, R 2 is NH 2 . Referring to formula (II), a subset of the purine compounds are those in which each of R 1 is NH 2 and R 2 is H. In these compounds, R 3 can be H; X can be O; and each of m and n, independently, is 1 or 2. In some embodiments, m is 2, X is O, n is 1, and Y is P(O)(OR a )(OR b ), in which R a can be H, and R b can be 5-ethylidene-(3,4-dialkoxy)-furan-2-one. Another subset of the purine compounds of formula (II) are those in which R 1 is OH and R 2 is NH 2 . In these compounds, R 3 can be H; X can be O; and each of m and n, independently, is 1 or 2. In some embodiments, m is 2, X is O, n is 1, and Y is P(O)(OR a )(OR b ), in which R a can be H, and R b can be 5-ethylidene-(3,4-dialkoxy)-furan-2-one. Also within the scope of this invention is a method of preparing certain purine compounds of formula (I). The method includes reacting a compound of formula (III): with an alkyl-X—(CH 2 ) halide to obtain a compound of formula (IV): In formulae (III) and (IV), R 1 is NH 2 or OH; R 2 is H or NH 2 ; R 3 is H or alkyl; and X is O, S, or NH; provided that if R 1 is NH 2 , R 2 is H; and if R 1 is OH, R 2 is NH 2 . Unless specifically pointed out, alkyl, aryl, heteroaryl, cyclyl, heterocyclyl, adenine, guanine, cytosine, uracil, and thymine mentioned above include both substituted and unsubstituted moieties. The term “substituted” refers to one or more substituents (which may be the same or different), each replacing a hydrogen atom. Examples of substituents include, but are not limited to, halogen, cyano, nitro, hydroxyl, amino, mercapto, C 1 ˜C 6 alkyl, C 1 ˜C 6 alkenyl, C 1 ˜C 6 alkynyl, aryl, heteroaryl, C 4 ˜C 8 cyclyl, C 4 ˜C 8 heterocyclyl, alkyloxy, aryloxy, alksulfanyl, arylsulfanyl, alkylamino, arylamino, dialkylamino, diarylamino, alkylcarbonyl, arylcarbonyl, heteroarylcarbonyl, alkylcarboxyl, arylcarboxyl, heteroarylcarboxyl, alkyloxycarbonyl, aryloxycarbonyl, heteroaryloxycarbonyl, alkylcarbamido, arylcarbamido, heterocarbamido, alkylcarbamyl, arylcarbamyl, heterocarbamyl, wherein each of alkyl (including alk), alkenyl, aryl, heteroaryl, cyclyl, and heterocyclyl is optionally substituted with halogen, cyano, nitro, hydroxyl, amino, mercapto, C 1 ˜C 6 alkyl, aryl, heteroaryl, alkyloxy, aryloxy, alkylcarbonyl, arylcarbonyl, alkylcarboxyl, arylcarboxyl, alkyloxycarbonyl, or aryloxycarbonyl. The term “alkyl” refers to both linear and branched alkyl. The term “cyclyl” refers to a hydrocarbon ring containing 4 to 8 carbons. The term “heterocyclyl” refers to a ring containing 4 to 8 ring members that have at least one heteroatom (e.g., S, N, or O) as part of the ring. The term “aryl” refers to a hydrocarbon ring system having at least one aromatic ring. Examples of aryl moieties include, but are not limited to, phenyl, naphthyl, and pyrenyl. The term “heteroaryl” refers to a hydrocarbon ring system having at least one aromatic ring which contains at least one heteroatom such as O, N, or S. Examples of heteroaryl moieties include, but are not limited to, furyl, pyrrolyl, thienyl, oxazolyl, imidazolyl, thiazolyl, pyridinyl, pyrimidinyl, quinazolinyl, and indolyl. All of the purine compounds described above include the compounds themselves, as well as their salts. The salts, for example, can be formed between a positively charged substituent (e.g., amino) on a compound and an anion. Suitable anions include, but are not limited to, chloride, bromide, iodide, sulfate, nitrate, phosphate, citrate, methanesulfonate, trifluoroacetate, and acetate. Likewise, a negatively charged substituent (e.g., carboxylate) on a compound can form a salt with a cation. Suitable cations include, but are not limited to, sodium ion, potassium ion, magnesium ion, calcium ion, and an ammonium cation such as teteramethylammonium ion. In addition, some of the just-described purine compounds have one or more double bonds, or one or more asymmetric centers. Such compounds can occur as racemates, racemic mixtures, single enantiomers, individual diastereomers, diastereomeric mixtures, and cis- or trans- or E- or Z- double bond isomeric forms. Exemplary purine compounds of formula (I) and (II) include: This invention also features a method of treating infection by virus. The method includes administering to a subject in need thereof an effective amount of a purine compound of formula (I) or formula (II) as described above. Examples of the viruses include DNA viruses such as herpesviridae (e.g., herpes simplex viruses) and retroviruses such as Moloney murine sarcoma viruses and human immunodeficiency viruses (e.g., human immunodeficiency viruses-1 or -2). As used herein, the term “treating infection” refers to use of one or more purine compounds described above for preventing or treating infection by virus, or other disease states secondary to viral infection, e.g., cervical cancer induced by Papilovirus. Also within the scope of this invention are a composition containing one or more of the aforementioned purine compounds for use in treating viral infection, and the use of such a composition for the manufacture of a medicament for infection treatment. Other features or advantages of the present invention will be apparent from the following detailed description of several embodiments, and also from the appending claims. DETAILED DESCRIPTION This invention relates to the purine compounds of formula (I) and formula (II) as described in the Summary section. As shown in Scheme 1a below, the purine compounds of formula (I) can be prepared by a novel procedure: alkylation of N 9 -tritylated purine compounds (III), followed by concomitant self-detritylation to yield the desired N 7 -alkylated purine nucleosides (IV). In this Scheme, R 1 , R 2 , R 3 , and X are defined in the Summary section. Compound (III) can be obtained by silylation of a purine (e.g., guanine) with hexamethyldisilazane in the presence of a catalytic amount of (NH 4 ) 2 SO 4 at an elevated temperature, followed by condensation of the resultant silylated purine with trityl chloride. Other than the just-described procedure, the purine compounds of formula (I) and formula (II) can be prepared by methods well known in the art. Scheme 1b shown below adepicts synthesis of the compounds of formula (I). In this scheme, R 1 , R 2 , R 3 , m, n, X, and Y are defined in the Summary section. More specifically, a purine compound is alkylated with 3-bromopropionitrile in the presence of NaH to give N 9 -cyanoethyl purine intermediate. Reaction of this intermediate with a methyliodo-ester and lithium 2,2,6,6-tetramethylpiperidine affords a mixture of N 7 -alkylated and N 9 -alkylated ester-containing products. Reduction of the ester in N 9 -alkylated product gives another intermediate, which is converted to a desired purine compound of formula (I) by an alkylation reaction or by reacting with a phosphonate in the presence of tert-butoxide. Similarly, reduction of the ester in N 7 -alkylated product gives a compound, which reacts with a phosphonate to produce an N 7 -substituted phosphonate-purine. Reaction of the N 7 -substituted phosphonate-purine with a halide in the presence of NaHCO 3 affords a desired compound of formula (II). Scheme 1c below depicts synthesis of the compounds of formula (II). In this scheme, R 1 , R 2 , R 3 , R a , R b , m, n, X, and Y are defined in the Summary section. The chemicals used in the above-described synthetic routes may include, for example, solvents, reagents, catalysts, protecting group and deprotecting group reagents. The methods described above may also additionally include steps, either before or after the steps described specifically herein, to add or remove suitable protecting groups in order to ultimately allow synthesis of the purine compound. In addition, various synthetic steps may be performed in an alternate sequence or order to give the desired compounds. Synthetic chemistry transformations and protecting group methodologies (protection and deprotection) useful in synthesizing applicable purine compounds are known in the art and include, for example, those described in R. Larock, Comprehensive Organic Transformations , VCH Publishers (1989); T. W. Greene and P. G. M. Wuts, Protective Groups in Organic Synthesis , 2 nd Ed., John Wiley and Sons (1991); L. Fieser and M. Fieser, Fieser and Fieser's Reagents for Organic Synthesis , John Wiley and Sons (1994); and L. Paquette, ed., Encyclopedia of Reagents for Organic Synthesis , John Wiley and Sons (1995) and subsequent editions thereof. A purine compound thus synthesized can be further purified by a method such as column chromatography, high pressure liquid chromatography, or recrystallization. Also within the scope of this invention is a pharmaceutical composition that contains an effective amount of at least one purine compound of the present invention and a pharmaceutically acceptable carrier. Further, this invention covers a method of administering to a subject in need of treating viral infection an effective amount of one or more of the purine compounds. The term “treating” is defined as the application or administration of a composition including the purine compound to a subject, who has a viral infection, a symptom of the infection, a disease or disorder secondary to the infection, or a predisposition toward the infection, with the purpose to cure, alleviate, relieve, remedy, or ameliorate the infection, the symptom of the infection, the disease or disorder secondary to the infection, or the predisposition toward the infection. “An effective amount” is defined as the amount of the purine compound which, upon administration to a subject in need thereof, is required to confer therapeutic effect on the subject. An effective amount of the purine compound may range from 5 mg/Kg to 20 mg/Kg. Effective doses also vary, as recognized by those skilled in the art, depending on route of administration, excipient usage, and the possibility of co-usage with any other therapeutic agent, such as an antiviral agent. To practice the method of the present invention, the above-described pharmaceutical composition can be administered orally, parenterally, by inhalation spray, topically, rectally, nasally, buccally, vaginally or via an implanted reservoir. The term “parenteral” as used herein includes subcutaneous, intracutaneous, intravenous, intramuscular, intraarticular, intraarterial, intrasynovial, intrastemal, intrathecal, intralesional and intracranial injection or infusion techniques. A sterile injectable composition, e.g., a sterile injectable aqueous or oleaginous suspension, can be formulated according to techniques known in the art using suitable dispersing or wetting agents (such as Tween 80) and suspending agents. The sterile injectable preparation can also be a sterile injectable solution or suspension in a non-toxic parenterally acceptable diluent or solvent, for example, as a solution in 1,3-butanediol. Among the acceptable vehicles and solvents that can be employed are mannitol, water, Ringer's solution and isotonic sodium chloride solution. In addition, sterile, fixed oils are conventionally employed as a solvent or suspending medium (e.g., synthetic mono- or di-glycerides). Fatty acids, such as oleic acid and its glyceride derivatives are useful in the preparation of injectables, as are natural pharmaceutically-acceptable oils, such as olive oil or castor oil, especially in their polyoxyethylated versions. These oil solutions or suspensions can also contain a long-chain alcohol diluent or dispersant, or carboxymethyl cellulose or similar dispersing agents. Other commonly used surfactants such as Tweens or Spans or other similar emulsifying agents or bioavailability enhancers which are commonly used in the manufacture of pharmaceutically acceptable solid, liquid, or other dosage forms can also be used for the purposes of formulation. A composition for oral administration can be any orally acceptable dosage form including, but not limited to, capsules, tablets, emulsions and aqueous suspensions, dispersions and solutions. In the case of tablets for oral use, carriers which are commonly used include lactose and corn starch. Lubricating agents, such as magnesium stearate, are also typically added. For oral administration in a capsule form, useful diluents include lactose and dried corn starch. When aqueous suspensions or emulsions are administered orally, the active ingredient can be suspended or dissolved in an oily phase combined with emulsifying or suspending agents. If desired, certain sweetening, flavoring, or coloring agents can be added. A nasal aerosol or inhalation composition can be prepared according to techniques well-known in the art of pharmaceutical formulation and can be prepared as solutions in saline, employing benzyl alcohol or other suitable preservatives, absorption promoters to enhance bioavailability, fluorocarbons, and/or other solubilizing or dispersing agents known in the art. A purine compound-containing composition can also be administered in the form of suppositories for rectal administration. The carrier in the pharmaceutical composition must be “acceptable” in the sense of being compatible with the active ingredient of the formulation (and preferably, capable of stabilizing it) and not deleterious to the subject to be treated. For example, solubilizing agents such as cyclodextrins, which form specific, more soluble complexes with the indole compounds, or one or more solubilizing agents, can be utilized as pharmaceutical excipients for delivery of the purine compounds. Examples of other carriers include colloidal silicon dioxide, magnesium stearate, cellulose, sodium lauryl sulfate, and D&C Yellow #10. A purine compound of this invention can be preliminarily screened for its efficacy in treating viral infection by one or more of the following in vitro assays. In one assay, a purine compound is tested for its inhibition of cytopathogenicity against the herpes simplex type 1 virus, herpes simplex simplex type 2 virus, thymidine kinase-positive and thymidine kinase-deficient strains of varicell-zoster virus, or human cytomegalovirus in Vero cells. The method for measuring viruses-induced cytophogenicity in Vero cell cultures, as well as the toxicity of the test compound toward HeLa and Vero cells, has been described in e.g., De Clercq et al. (1980) J. Infect. Dis . 141: 563-574. In another assay, a purine compound is tested for its inhibition of cytopathogenicity against the human immunodeficiency viruses HIV-1 (IIIB) and HIV-2 (LAV-2) in MT4 cells. The method for measuring viruses-induced cytophogenicity in MT4 cells or CEM cells, as well as the toxicity toward MT4 and CEM cells, has been described in e.g., Averett, (1989) J. Virol. Methods 23: 263-276. The antiviral activity of a purine compound can be further assessed using an in vivo animal model. See the specific examples below. Without further elaboration, it is believed that the above description has adequately enabled the present invention. The following specific embodiments are, therefore, to be construed as merely illustrative, and not limitative of the remainder of the disclosure in any way whatsoever. All of the publications cited herein are hereby incorporated by reference in their entirety. General. For anhydrous reactions, glassware was dried overnight in an oven at 120° C. and cooled in a desiccator over anhydrous CaSO 4 or silica gel. Reagents were purchased from Fluka and enzymes from Sigma Chemical Co. Solvents, including dry ether and tetrahydrofuran (THF), were obtained by distillation from the sodium ketyl of benzophenone under nitrogen. Other solvents, including chloroform, dichloromethane, ethyl acetate, and hexanes were distilled over CaH 2 under nitrogen. Absolute methanol and ethanol were purchased from Merck and used as received. Melting points were obtained with a Büchi 510 melting point apparatus. Infrared (IR) spectra were recorded on a Beckman IR-8 spectrophotometer. The wavenumbers reported are referenced to the 1601 cm −1 absorption of polystyrene. Proton NMR spectra were obtained on a Varian XL-300 (300 MHz) Spectrometer. Chloroform-d and dimethylsulfoxide-d 6 were used as solvent; Me 4 Si (δ 0.00 ppm) was used as an internal standard. All NMR chemical shifts are reported as 6 values in parts per million (ppm) and coupling constants (J) are given in hertz (Hz). The splitting pattern abbreviations are as follows: s, singlet; d, doublet; t, triplet; q, quartet; br, broad; m, unresolved multiplet due to the field strength of the instrument; and dd, doublet of doublets. UV spectroscopy was carried out using an HP8452A diode array spectrophotometer. Mass spectra were carried out on a VG 70-250 S mass spectrometer. Microanalyses were performed on a Perkin-Elmer 240-B microanalyzer. Purification on silica gel refers to gravity column chromatography on Merck Silica Gel 60 (particle size 230-400 mesh). Analytical TLC was performed on precoated plates purchased from Merck (Silica Gel 60 F 254 ). Compounds were visualized by use of UV light, I 2 vapor, or 2.5% phosphomolybdic acid in ethanol with heating. 9-[(2-Hydroxyethoxy)methyl]adenine 1,9-[(2-hydroxyethoxy)methyl]guanine (acyclovir 2), 9-(β-D-arabinofaranosyl)adenine (ara-A 3), 9-[2-(phosphonomethoxy)ethyl]adenine (PMEA 4) were used for comparison in Examples 8-10. Compounds 9, 14, 20, 24, 25, 28, and 29 were prepared by the methods described in Examples 1-7. EXAMPLE 1 Synthesis of Compound 9 As shown in Scheme 2, N 9 -tritylated guanine 6 was synthesized in two steps. Silylation of guanine 5 with hexamethyldisilazane (HMDS) in the presence of a catalytic amount of (NH 4 ) 2 SO 4 at refluxing temperature followed by condensation of the resultant silylated guanine with trityl chloride in MeCN at 25° C. afforded the desired N 9 -tritylated guanine 6 in 90% yield. Treatment of 6 with (chloroethoxy)methyl chloride (Hakimelahi & Khalafi-Nezhad (1989) Helv. Chim. Acta 72: 1495-1550) in DMF at room temperature gave the corresponding N 7 -alkylated guanine 7 in 88% yield. Likewise, treatment of 6 with [2-(p-methoxyphenyloxy)ethoxy]methyl chloride (Khorshidi (1986) Doctoral Thesis in Pharmacy, Faculty of Medicine, Isfahan University, Isfahan, Iran) led to the N 7 -isomer 8 in 82% yield. Removal of the p-methoxyphenyl moiety was then achieved by treatment with ceric ammonium nitrate (CAN, Fukuyama et al. (1985) Tetrahedron Lett . 26: 6291-6292) in a mixture of MeCN and H 2 O (3:1) at 0-25° C. to afford compound 9 in 70% yield. For compound 9 (i.e., a N 7 -isomer) and its corresponding N 9 -isomer, their 1 H and 13 C NMR spectra are different. See, e.g., Kjellberg & Johansson (1986) Tetrahedron 42: 6541-6544; Shiragami et al. (1995) Nucleosides & Nucleotides 14: 337-340; and Bailey & Hamden (1987) Nucleosides & Nucleotides 6: 555-574. The 1 H signals of H 2 C(1′) (5.81 ppm) and HC(8) (8.67 ppm) for the N 7 -isomer were found to be shifted downfield relative to the corresponding signals of the N 9 -isomer, in which the H 2 C(1′) and HC(8) resonated at 5.35 and 7.81 ppm, respectively. On the other hand, the 1 H signals for NH 2 was observed to be shifted upfield for the N 7 -isomer, 5.96 ppm for the N 7 -isomer, relative to the corresponding signal for N 9 -isomer, which was observed at 6.52 ppm. The 13 C NMR signals for C(1′) (75.25 ppm) and C(8) (143.92 ppm) of the N 7 -isomer were found to be shifted downfield relative to the corresponding signals of the N 9 -isomer, which were observed respectively at 71.64 and 137.89 ppm. In contrast, the signal of C(5) of the N 7 -isomer resonated at 107.16 ppm, which was upper field to that of the N 9 -isomer at 116.52 ppm. The UV λ max of the N 7 -isomer appeared at 289 nm, whereas the corresponding λ max of the N 9 -isomer appeared at 253 and 273 (sh) nm. 9-(Triphenylmethyl)guanine 6. Guanine 5 (1.51 g, 9.99 mmol) and (NH 4 ) 2 SO 4 (100 mg) were suspended in HMDS (150 mL) and refluxed for 24 h. The solvent was evaporated under reduced pressure and the residue was dissolved in CH 3 CN (150 mL). Triphenylmethyl chloride (2.79 g, 10.0 mmol) was added and the reaction mixture was stirred at 25° C. for 7.0 h. The solution was concentrated under reduced pressure and the residue was purified by use of column chromatography (hexanes/EtOAc=1.5:8.5) to afford 6 (3.54 g, 8.99 mmol) in 90% yield: mp 268-270° C.; R f (hexanes/EtOAc=1:2) 0.34; UV (EtOH) λ max 254 (ε 13,870), 278 (sh); 1 H NMR (DMSO-d 6 ) δ 5.97 (s, 2 H, NH 2 ), 7.09-7.39 (m, 16 H, HC 8 +C(C 6 H 5 ) 3 ), 10.45 (br s, 1H, NH); MS m/z 393 (M + ). Anal. (C 24 H 19 N 5 O) C, H, N; calcd (%): 73.26, 4.87, 17.80; found (%): 73.20, 4.81, 17.78. 7-[(2-Chloroethoxy)methy]guanine 7. To a solution of 6 (1.77 g, 4.49 mmol) in DMF (30 mL) was added (2-chloroethoxy)methyl chloride (0.65 g, 5.0 mmol). The reaction mixture was stirred at 25° C. for 8.0 h. The solution was then partitioned between EtOAc (100 mL) and water (100 mL). The EtOAc solution was washed with water (4×100 mL); then it was dried over MgSO 4 (s) and filtered. Evaporation under reduced pressure and purification of the residue by use of column chromatography (EtOAc) afforded 7 (1.07 g, 4.39 mmol) in 88%: mp>280° C. (dec.); R f (hexanes/EtOAc=1:2) 0.20; UV (EtOH) λ max 288 (ε 15,100); 1 H NMR (DMSO-d 6 ) δ 3.67 (t, J=6.10 Hz, 2 H, CH 2 Cl), 3.75 (t, J=6.10 Hz, 2 H, OCH 2 ), 5.62 (s, 2 H, H 2 C 1′ ), 6.15 (s, 2 H, NH 2 ), 7.58 (s, 1 H, NH), 8.16 (s, 1 H, HC 8 ); 13 C NMR (DMSO-d 6 ) δ 43.39 (CH 2 Cl), 68.50 (OCH 2 ) 74.71 (C 1′ ), 107.71 (C 5 ), 143.87 (C 8 ), 153.15 (C 4 ), 154.37 (C 2 ), 159.12 (C 6 ); MS m/z 243 (M + , Cl-cluster). Anal (C 8 H 10 ClN 5 O 2 ) C, H, N, Cl; calcd (%): 39.43, 4.14, 28.74, 14.56; found (%): 39.38, 4.13, 28.70, 14.45. 7-[(2-(p-Methoxyphenyloxy)ethoxy)methyl]guanine 8. Compound 8 (2.72 g, 8.20 mmol) was prepared in 82% yield from 6 (3.64 g, 9.25 mmol) and (2-(p-methoxyphenyloxy)ethoxy)methyl chloride (2.18 g, 10.1 mmol) in DMF (100 mL) by the method used for the synthesis of 7: mp>250° C. (dec.); R f (hexanes/EtOAc=1:2) 0.18; UV (EtOH) λ max 290 (ε 16,500); 1 H NMR (DMSO-d 6 ) δ 3.70 (s, 3 H, OCH 3 ), 4.02-4.10 (m, 4 H, O(CH 2 ) 2 O), 5.74 (s, 2 H, H 2 C 1′ ), 5.88 (s, 2 H, NH 2 ), 6.67, 6.70 (AA′BB′, J=9.30 Hz, 4 H, C 6 H 4 ), 7.50 (s, 1 H, NH), 8.85 (s, 1 H, HC 8 ); 13 C NMR (DMSO-d 6 ) δ 56.12 (CH 3 ), 68.74 (OCH 2 ), 70.20 (CH 2 Oph), 75.31(C 1′ ), 107.25 (C 5 ), 144.02 (C 8 ), 153.25 (C 4 ), 154.48 (C 2 ), 159.26 (C 6 ), 115.55, 116.38, 153.51, 155.56 (C 6 H 4 ); MS m/z 331 (M + ). Anal (C 15 H 17 N 5 O 4 ) C, H, N; calcd (%): 54.37, 5.17, 21.13; found (%): 54.26, 5.20, 21.19. 7-[(2-Hydroxyethoxy)methyl]guanine 9. To a solution of 8 (1.36 g, 4.10 mmol) in a mixture of CH 3 CN (30 mL) and water (10 mL) was added CAN (2.25 g, 4.10 mmol) at 0° C. The stirred reaction mixture was allowed to warm-up to 25° C. within 1.0 h. Water (30 mL) was added to afford a solid. Filtration and crystallization of the solid from EtOH/water (4:1) gave 9 (0.65 g, 2.88 mmol) in 70% yield: mp>280° C. (dec.); R f (hexanes/EtOAc=1:2) 0.08; UV (EtOH) λ max 289 (ε 16,000); 1 H NMR (DMSO-d 6 ) δ 3.61-3.82 (m, 4H, O(CH 2 ) 2 O), 4.85 (br s, 1H, OH), 5.81 (s, 2 H, H 2 C 1′ ), 5.96 (br s, 2 H, NH 2 ), 6.61 (s, 1 H, NH), 8.67 (s, 1 H, HC 8 ); 13 C NMR (DMSO-d 6 ) δ 60.91 (CH 2 OH), 70.01 (OCH 2 ) 75.25 (C 1′ ), 107.16 (C 5 ), 143.92 (C 8 ), 153.20 (C 4 ), 154.40 (C 2 ), 159.30 (C 6 ); MS m/z 225 (M + ). Anal (C 8 H 11 N 5 O 3 ) C, H, N; calcd (%): 42.66, 4.92, 31.10; found (%): 42.75, 4.84, 31.21. EXAMPLE 2 Synthesis of Compound 14 As shown in Scheme 3, by the same synthetic strategy shown in Scheme 2, novel N 7 -alkylated adenines 12 and 14 were obtained from adenine 10 via the N 9 -tritylated adenine 11. Reaction of 11 with (2-chloroethoxy)methyl chloride in CH 2 Cl 2 at 25° C. gave N 7 -adenine derivative 12 in 90% yield. Likewise, reaction of 11 with [2-(p-methoxyphenyloxy)ethoxy]methyl chloride in CH 2 Cl 2 at 25° C. gave N 7 -isomer 13 in 86% yield. Treatment of 13 with CAN then produced the deprotected compound 14 in 80% yield. For compound 14 (i.e., a N 7 -isomer) and its corresponding N 9 -isomer, their 1 H and 13 C NMR spectra are different. In methanol-d 4 , pronounced downfield shifts were observed for the signals resulting from the H 2 C(1′) (5.79 ppm) and HC(8) (8.60 ppm) of the N 7 -isomer when compared with those of the N 9 -isomer, respectively observed at 5.67 and 8.27 ppm. In DMSO-d 6 , the NH 2 signals of the N 7 -isomer was shifted downfield to ˜9.2 and 110.0 ppm, whereas the corresponding signals of the N 9 -isomer appeared as a broad peak at ˜7.3 ppm. The 13 C NMR signals for C(1′) (82.01 ppm) and C(8) (147.68 ppm) for the N 7 -isomer were found to be shifted downfield relative to those of the N 9 -isomer which were observed respectively at 74.31 and 143.14 ppm. The signals of the C(2) (144.29 ppm) and C(6) (152.12 ppm) for the N 7 -isomer, however, were found to be shifted upfield relative to those of the N 9 -isomer, observed respectively at 153.74 and 157.19 ppm. Furthermore, the H—C(2) coupling constants for the N 7 - and N 9 -isomers were respectively 220 and 204 Hz, whereas the corresponding H—C(8) coupling constants were respectively 224 and 217 Hz. The attachment of the side chain at the N-7 position of adenine was confirmed by HMQC spectroscopy, in which the H 2 C(1′) and C(6) exhibited a strong interaction in the N 7 -isomer whereas the corresponding N 9 -isomer showed long-range coupling between H 2 C(1′) and C(4). 9-(Triphenylmethyl)adenine 11. To a solution of 10 (1.35 g, 9.99 mmol) and DMF (10 mL) in pyridine (30 mL) was added Ph 3 Cl (2.79 g, 10.0 mmol). The reaction mixture was stirred at 25° C. for 7.0 h. It was then diluted with EtOAc (150 mL) and water (100 mL). The organic layer was separated and then washed with water (4×100 mL). It was dried over MgSO 4 (s) and concentrated under reduced pressure to yield a foam. Purification was carried out by use of column chromatography (EtOAc/hexanes=8:2) to afford compound 11 (3.58 g, 9.48 mmol) in 95% yield: mp 260-262° C.; R f (hexanes/EtOAc=1:2) 0.38; UV (EtOH) λ max 260 (ε 14,600); 1 H NMR (CDCl 3 ) δ 5.60 (s, 2 H, NH 2 ), 7.15-7.31 (m, 15 H, C(C 6 H 5 ) 3 ), 7.74 (s, 1 H, HC 2 ), 8.05 (s, 1 H, HC 8 ); MS m/z 377 (M + ); Anal. (C 24 H 19 N 5 ) C, H, N; calcd. (%): 76.37, 5.07, 18.55; found (%): 76.39, 5.16, 18.60. 7-[(2-Chloroethoxy)methy]adenine 12. To a solution of 11 (1.79 g, 4.74 mmol) in CH 2 Cl 2 (70 mL) was added (2-chloroethoxy)methyl chloride (0.65 g, 5.0 mmol). The reaction mixture was stirred at 25° C. for 8.0 h to afford a solid. Filteration and crystallization from MeOH gave 12 (1.66 g, 4.27 mmol) in 90% yield: mp 184-186° C.; R f (hexanes/EtOAc=1:2) 0.28; UV (EtOH) λ max 267 (ε 14,600); 1 H NMR (DMSO-d 6 ) δ 3.73 (t, J=4.5 Hz, 2 H, CH 2 Cl), 3.86 (t, J=4.5 Hz, 2 H, OCH 2 ), 5.76 (s, 2 H, H 2 C 1′ ), 8.49 (s, 1 H, HC 2 ), 8.74 (s, 1 H, HC 8 ), 9.32, 10.05 (2 br, 2 H, NH 2 ); 13 C NMR (DMSO-d 6 ) δ 43.13 (CH 2 Cl), 68.82 OCH 2 78.10 (C 1′ ), 118.20 (C 5 ), 143.53 (C 2 ), 146.50 (C 8 ), 149.39 (C 4 ), 151.60 (C 6 ); MS m/z 227 (M + , Cl-cluster). Anal. (C 8 H 10 ClN 5 O) C, H, N, Cl; calcd (%): 42.20, 4.43, 30.76, 15.59; found (%): 42.32, 4.50, 30.69, 15.62. 7-[(2-(p-Methoxyphenyloxy)ethoxy)methyl]adenine 13. After 10 h, compound 13 (2.79 g, 8.86 mmol) was synthesized in 86% yield from 11 (3.90 g, 10.3 mmol) and (2-(p-methoxyphenyloxy)ethoxy)methyl chloride (2.25 g, 10.4 mmol) in CH 2 Cl 2 (100 mL) by the method used for the synthesis of 12: mp 129-130° C.; R f (hexanes/EtOAc=1:2) 0.21; UV (EtOH) λ max 270 (ε 15,100); 1 H NMR (CD 3 OD) δ 3.68 (s, 3 H, OCH 3 ), 4.04 (br s, 4 H, O(CH 2 ) 2 O), 5.83 (s, 2 H, H 2 C 1′ ), 6.58, 6.68 (AA′BB′, J=9.0 Hz, 4 H, C 6 H 4 ), 8.29 (s, 1 H, HC 2 ), 8.61 (s,1 H, HC 8 ); 13 C NMR (CD 3 OD) δ 56.03 (CH 3 ), 68.70 (OCH 2 ), 70.18 CH 2 OPh 82.27 (C 1′ ), 119.59 (C 5 ), 144.17 (C 2 ), 147.57 (C 8 ), 150.25 (C 4 ), 151.97 (C 6 ), 115.50, 116.35, 153.44, 155.53 (C 6 H 4 ); MS m/z 315 (M + ). Anal. (C 15 H 17 N 5 O 3 ) C, H, N; calcd (%): 57.13, 5.43, 22.21; found (%): 57.24, 5.49, 22.34. 7-[(2-Hydroxyethoxy)methyl]adenine 14. Compound 14 (0.74 g, 3.54 mmol) was prepared in 80% yield from 13 (1.39 g, 4.43 mmol) and CAN (2.43 g, 4.43 mmol) in CH 3 CN (30 mL) and water (10 mL) by the method used for the synthesis of 9: mp 122-123° C.; R f (hexanes/EtOAc=1:2) 0.12; UV (EtOH) λ max 265 (ε 14,800); 1 H NMR (CD 3 OD) δ 3.69-3.79 (m, 4 H, O(CH 2 ) 2 O), 5.79 (s, 2 H, H 2 C 1′ ), 8.32 (s, 1 H, HC 2 ), 8.60 (s, 1 H, HC 8 ); 13 C NMR (CD 3 OD) δ 61.80 (CH 2 OH), 72.27 (OCH 2 ) 82.01 (C 1′ ), 119.62 (C 5 ), 144.29 (C 2 ), 147.68 (C 8 ), 150.32 (C 4 ), 152.12 (C 6 ); MS m/z 209 (M + ). Anal (C 8 H 11 N 5 O 2 ) C, H, N; calcd (%): 45.93, 5.30, 33.48; found (%): 45.80, 5.45, 33.51. EXAMPLE 3 Synthesis of Compound 20 As shown in Scheme 4, alkylation of adenine 10 with 3-bromopropionitrile in the presence of NaH in DMF gave N 9 -(cyanoethyl)adenine 15 in 75% yield. Reaction of 15 with methyl iodoacetate and lithium 2,2,6,6-tetramethylpiperidine (Li TMP) in THF afforded a mixture of N 7 -alkylated product 16 (60% yield) and N 9 -isomer 17 (20% yield). Reduction of the ester group in 16 with NaBH 4 in wet THF 13 gave N 7 -(hydroxyethyl)adenine 18 in 55% yield. Conversion of 18 to phosphonate 19 (60% yield) was accomplished by use of diethyl (p-toluenesulfonyloxymethane)phosphonate and sodium tert-butoxide in DMF. 14 Treatment of compound 19 with Me 3 SiBr 15 then afforded phosphonic acid 20 in 45% yield. 9-(2-Cyanoethyl)adenine 15. To a suspension of 57% NaH in mineral oil (0.962 g, 22.8 mmol) in dry DMF (100 mL) was added 10 (3.40 g, 24.9 mmol) under nitrogen and the mixture was heated at 80° C. for 1.0 h. A solution of 3-bromopropionitrile (2.81 g, 21.0 mmol) in DMF (5.0 mL) was added at 25° C. and the reaction was heated at 65° C. for 17 h. It was then diluted with EtOAc (250 mL) and 5% aqueous HCl solution (150 mL). The organic layer was separated and then washed with water (4×100 mL). It was dried over MgSO 4 (s) and concentrated under reduced pressure. Purification of the residue was carried out by use of column chromatography (EtOAc/hexanes=7.5:2.5) to afford compound 15 (3.510 g, 18.67 mmol) in 75% yield: mp 148-150° C.; R f (hexanes/EtOAc=1:2) 0.38; UV (EtOH) λ max 260 (ε 13,900); 1 H NMR (CD 3 OD) δ 2.92 (t J=5.8 Hz, 2 H, CH 2 CN),3.43 (t, J=5.8 Hz, 2 H, CH 2 N), 8.35 (s, 1 H, HC 2 ), 8.69 (s, 1 H, HC 8 ); MS m/z 188 (M + ). 7-[(Methoxycarbonyl)methyl]adenine 16 and 9-[(Methoxycarbonyl)methyl]-adenine 17. To a stirred solution containing 15 (0.134 g, 1.00 mmol) and methyl iodoacetate (0.30 g, 1.5 mmol) in dry THF (20 mL) was added a THF solution of LiTMP (2.8 mL, 1.2 mmol) dropwise under an argon atmosphere at −20° C. The reaction mixture was warmed to 25° C. within 1.0 h; then it was stirred at room temperature for 15 h. The solution was partitioned between EtOAc (40 mL) and water (40 mL). The organic layer was dried over MgSO 4 (s) and concentrated under reduced pressure. Purification of the residue by use of column chromatography (EtOAc/hexanes=8:2) gave compound 17 (0.041 g, 0.20 mmol) in 20% yield. Further elution of the column with EtOAc/hexanes (9:1) afforded 16 (0.123 g, 0.594 mmol) in 60% yield. For 16: mp 112-114° C.; R f (hexanes/EtOAc=1:2) 0.16; UV (EtOH) λ max 267 (ε 14,000); 1 H NMR (CD 3 OD) δ 3.83 (s 3H, CH 3 ), 4.37 (s, 2 H, H 2 C 1′ ), 8.41 (s, 1 H, HC 2 ), 8.75 (s, 1 H, HC 8 ); MS m/z 207 (M + ). Anal (C 8 H 9 N 5 O 2 ) C, H, N; calcd (%): 46.37, 4.38 33.80; found (%): 46.32, 4.35, 33.83. For 17: mp 148-149° C.; R f (hexanes/EtOAc=1:2) 0.29; UV (EtOH) λ max 260 (ε 13,750); 1 H NMR (CD 3 OD) δ 3.75 (s, 3H, CH 3 ), 4.15 (s, 2 H, H 2 C 1′ ), 8.29 (s, 1 H, HC 2 ), 8.38 (s, 1 H, HC 8 ); MS m/z 207 (M + ). 7-(2-Hydroxyethyl)adenine 18. To a stirred solution containing 16 (0.207 g, 0.999 mmol) and water (0.50 mL) in THF (12 mL) was added NaBH 4 (0.38 g, 10.0 mmol). After stirring for 4 h at 25° C., the reaction mixture was neutralized to pH=7.0 by use of 10% HCl aqueous solution. Solvent was evaporated under reduced pressure and the residue was purified by column chromatography (EtOAc/hexanes=9:1) to give 18 (0.10 g, 0.55 mmol) in 55% yield: mp 132-133° C.; R f (hexanes/EtOAc=1:2) 0.09; UV (EtOH) λ max 265 (ε 14,110); 1 H NMR (CD 3 OD) δ 3.68 (t, J=5.9 Hz, 2 H, CH 2 O), 3.66 (t, J=5.9 Hz, 2 H, CH 2 N), 8.34 (s, 1 H, HC 2 ), 8.68 (s, 1 H, HC 8 ); MS m/z 179 (M + ). Anal (C 7 H 9 N 5 O) C, H, N; calcd (%): 46.92, 5.06, 39.08; found (%): 47.01, 5.12, 39.11. 7-[2-(D)iethylphosphonomethoxy)ethyl]adenine 19. To a solution of 18 (0.18 g, 0.10 nmmol) in DMF (15 mL) was added sodium tert-butoxide (0.150 g, 1.56 mmol). After 5 min, diethyl (p-toluenesulfonyloxymethane)phosphonate (0.42 g, 1.3 mmol) was added and the reaction mixture stirred at 35° C. for 4.0 h. The reaction was then quenched with acetic acid (5.0 mL) and the mixture was partitioned between EtOAc (50 mL) and water (50 mL). The organic layer was separated and then washed with water (5×60 mL), dried over MgSO 4 (s), filtered, and concentrated under reduced pressure. Purification by use of column chromatography (EtOAc/MeOH=9:1) gave 19 (0.20 g, 0.60 mmol) in 60% yield as a white foam: R f (hexanes/EtOAc=1:2) 0.16; UV (EtOH) λ max 266 (ε 15,321); 1 H NMR (CD 3 OD) δ 1.39 (t, J=6.7 Hz, 6H, 2 CH 3 ), 3.74 (t, J=6.3 Hz, 2 H, CH 2 N), 3.76 (d, J=9.0 Hz, 2 H, CH 2 P), 3.97-4.29 (m, 6 H, CH 2 O+2 CH 2 OP), 8.37 (s, 1 H, HC 2 ), 8.77 (s, 1 H, HC 8 ); MS m/z 329 (M + ). Anal (C 12 H 20 N 5 O 4 P) C, H, N; calcd (%): 43.76, 6.12, 21.27; found (%): 43.87, 6.21, 21.36. 7-[2-(Phosphonomethoxy)ethyl]adenine 20 To a solution of 19 (3.29 g, 10.0 mmol) in CH 2 Cl 2 (130 mL) and DMF (10 mL) was added Me 3 SiBr (4.95 g, 30.0 mmol); then the solution was stirred at 25° C. for 7.0 h. A mixture of MeOH and water (5:1, 40 mL) was added, and solvents were evaporated. Purification by use of column chromatography (EtOAc/MeOH=6:4) afforded 20 (1.23 g, 4.50 mmol) in 45% yield: mp 253° C. (dec.); R f (hexanes/EtOAc=1:2) 0.05; UV (EtOH) λ max 265 (ε 14,700); 1 H NMR (CD 3 OD) δ 3.66 (t, J=6.4 Hz, 2 H, CH 2 N), 3.71 (d, J=8.7 Hz, 2 H, CH 2 P), 4.19 (t, J=6.4 Hz, 2 H, CH 2 O), 8.36 (s, 1 H, HC 2 ), 8.78 (s, 1 H, HC 8 ); MS m/z 273 (M + ). Anal (C 8 H 12 N 5 O 4 P) C, H, N; calcd (%): 35.17, 4.43, 25.63; found (%): 35.21, 4.41, 25.71. EXAMPLES 4-7 Synthesis of Compound 24, 25, 28, and 29 As shown in Schemes 5 and 6, the nucleotide analog 24 was readily obtained in three steps from adenosine 5′-monophosphate 21, starting with silylation of phosphate 21 in CH 3 CN with t-BuMe 2 SiCl in the presence of AgNO 3 and pyridine. See, e.g., Ogilvie et al. (1983) Can. J. Chem . 61: 1204-1212. The resulting trisilylated compound 22 was condensed with phosphonic acid 20 using trichloromethanesulfonyl chloride in collidine and THF to afford dinucleotide 5′-monophosphate 23 in 42% overall yield. See, e.g., Hakimelahi et al. (1995) J. Med. Chem . 38: 4648-4659 and references cited therein. Desilylation of 23 with n-Bu 4 NF in THF at 25° C. gave dinucleotide 24 in 90% yield which was then reacted with (Z)-4-(2-chloroethylidene)-2,3-dimethoxy-Δ a,b -butenolide (26) in the presence of NaHCO 3 in DMF to afford the target molecule 25 in 85% yield (Hakimelahi et al. (2001) J. Med. Chem . 44: 1749-1757). Similarly, treatment of compound 26 with either 9-[2-(phosphonomethoxy)ethyl]adenine (PMEA 4) or 9-[2-(phosphonomethoxy)ethyl]guanine (PMEG 27) in the presence of NaHCO 3 in DMF respectively gave an 80% or 88% yield of the desired compound 28 or 29. 9-[2′-O-(tert-butyldimethylsilyl)-5′-O-(phosphono)-β-D-furanosyl]adenine-3′-[[1-(adenin-7-yl-ethoxy)methyl]phosphonate] 23. To a solution of adenosine 5′-monophosphate (21) monohydrate (3.65 g, 9.99 mmol) in a mixture of pyridine (150 mL) and CH 3 CN (160 mL) was added AgNO 3 (6.63 g, 39.0 mmol). After 10 min, tert-butyldimethylsilyl chloride (5.70 g, 37.8 mmol) was added. The mixture was stirred at 25° C. for 7.0 h and then filtered to remove AgCl. The filtrate was evaporated and the resultant crude product 22 was dissolved in dry THF (40 mL). In another flask, collidine (3.66 g, 30.0 mmol) was added to a solution of THF (45 mL) containing 20 (2.73 g, 10.0 mmol) at −10° C. To this solution was added CCl 3 SO 2 Cl (2.20 g, 10.0 mmol) in THF (15 mL) dropwise. After crude 22 in THF was added to the mixture, it was stirred at 25° C. for 10 h. The solvents were removed, and the residue was dissolved in AcOEt (100 mL) and washed with water (3×100 mL). The organic layer was concentrated, and the residue was purified by use of column chromatography (EtOAc/MeOH=6:4) to afford 23 (3.0 g, 4.2 mmol) in 42% overall yield: mp 223-225° C.; UV (EtOH) λ max 264 (ε 17,300); 1 H NMR (CD 3 OD) δ 0.16 (brs, 6 H, (CH 3 ) 2 Si), 1.05 (s, 9 H, (CH 3 ) 3 C), 3.67-4.27 (m, 8 H, CH 2 N+CH 2 O+CH 2 OP+CH 2 P), 4.32-4.5 (m, 3H, HC 2 ′+HC 3 ′+HC 4 ′), 6.58 (d, J=4.9 Hz, 1 H, HC 1 ′), 8.12, 8.42 (2s, 2 H, 2×HC 2 ), 8.27, 8.89 (2s, 2 H, 2×HC 8 ). Anal (C 24 H 38 N 10 P 2 Si) C, H, N; calcd (%) 40.22, 5.34, 15.94; found (%): 40.29, 5.38, 19.50. 9-[5′-O-(phosphono)-β-D-furanosyl]adenine-3′-[[1-(adenin-7-yl-ethoxy)methyl]phosphonate] 24. To a solution of 23 (0.36 g, 0.50 mmol) in THF (5.0 mL) was added n-Bu 4 NF (1.0 M solution in THF, 0.31 g, 1.2 mmol). Acetic acid (0.50 mL) was added to the mixture after it was stirred at 25° C. for 30 min. The solvents were removed, and the residue was purified by use of Whatman 3-mm paper with a mixture of i-PrOH, NH 4 OH, and H 2 O (9:1:2) as the eluent. The band at ca. R f 0.35 was eluted with H 2 O and collected by lyophilization to give 24 (0.27 g, 0.45 mmol) in 90% yield: mp>250° C. dec; UV (EtOH) λ max 264 (ε 18,200); 1 H NMR (CD 3 OD) δ 3.75-4.18 (m, 8 H, CH 2 N+CH 2 O+CH 2 OP+CH 2 P), 4.29-4.70 (m, 3H, HC 2 ′+HC 3 ′+HC 4 ′), 6.48 (d, J=4.5 Hz, 1 H, HC 1 ′), 7.99, 8.39 (2s, 2 H, 2×HC 2 ), 8.26, 8.83 (2s, 2 H, 2×HC 8 ). Anal (C 18 H 24 N 10 O 10 P 2 ) C, H, N; calcd (%): 35.88, 4.02, 23.25; found (%): 35.82, 4.12, 23.17. 9-[[(Z)-4-(ethylidene)-2,3-dimethoxy-Δ a,b -butenolide]-β-D-furanosyl]adenine-3′-[[1-(adenin-7-yl-ethoxy)methyl]phosphonate]-4,5′-phosphate 25. To a solution of 24 (0.300 g, 0.499 mmol) in DMF (20 mL) was added NaHCO 3 (0.30 g, 3.6 mmol). The reaction mixture was stirred at 25° C. under N 2 for 10 min. Then, butenolide 26 (0.10 g, 0.50 mmol) was added and stirred under N 2 for 1.0 h. The solution was diluted with EtOAc (50 mL) and aqueous HCl solution (1%, 40 mL). The organic layer was separated and washed with H 2 O (50 mL). Then, it was dried over MgSO 4 (s), filtered, and concentrated under reduced pressure. Purification by use of silica gel column chromatography with EtOAc/MeOH (6:4) as eluant afforded 25 (0.32 g, 0.42 mmol) in 85% yield: mp>237° C. dec; UV (EtOHl) λ max 215 (ε 16,000), λ max 264 (ε 18,540); 1 H NMR (CD 3 OD) δ 3.69-4.12 (m, 16 H, CH 2 N+CH 2 O+2×CH 2 OP+CH 2 P+C 2 OCH 3 +C 3 OCH 3 ), 4.31-4.78 (m, 3H, HC 2 ′+HC 3 ′+HC 4 ′), 5.38 (t, J=7.0 Hz, 1 H, ═CH), 6.51 (d, J=4.8 Hz, 1 H, HC 1 ′), 8.02, 8.40 (2s, 2 H. 2×HC 2 ), 8.28, 8.86 (2s, 2 H, 2×HC 8 ). Anal (C 26 H 32 N 10 O 14 P 2 ) C, H, N; calcd (%): 40.52, 4.19, 18.18; found (%): 40.61, 4.22, 18.21. [1-(Adenin-9-yl-ethoxy)methyl]phosphono-6-yl-(Z)-4-(ethylidene)-2,3-dimethoxy-Δ a,b -butenolide 28. Compound 28 (3.90 g, 8.80 mmol) was prepared in 88% yield from 4 (2.73 g, 9.99 mmol) and 26 (2.20 g, 10.0 mmol) in the same manner that 25 was prepared from 24: mp>241° C. (dec.); R f (hexanes/EtOAc=1:2) 0.12; UV (EtOH) λ max 218 (ε 13,097), λ max 259 (ε 14,700); 1 H NMR (CD 3 OD) δ 3.57 (t, J=6.0 Hz, 2 H, CH 2 N), 3.69 (d, J=9.0 Hz, 2 H, CH 2 P), 3.89 (m, 5 H, C 2 OCH 3 +CH 2 OP), 4.06, (t, J=6.0 Hz, 2 H, CH 2 O), 4.13 (s, 3 H, C 3 OCH 3 ), 5.41 (t, J=7.0 Hz, 1H, ═CH), 8.12 (s, 1 H, HC 2 ), 8.21 (s, 1 H, HC 8 ); MS m/z 441 (M + ). Anal (C 16 H 20 N 5 O 8 P) C, H, N; calcd (%): 43.54, 4.57, 15.86; found (%): 43.66, 4.46, 15.95. [1-(Guanine-9-yl-ethoxy)methyl]phosphono-6-yl-(Z)-4-(ethylidene)-2,3-dimethoxy-Δ a,b -butenolide 29. Compound 29 (3.7 g, 8.0 mmol) was prepared in 80% yield from 27 (2.89 g, 9.99 mmol) and 26 (2.20 g, 10.0 mmol) in the same manner that 25 was prepared from 24: mp>260° C. (dec.); R f (hexanes/EtOAc=1:2) 0.05; UV (EtOH) λ max 252 (ε 11,097), λ max 273 (ε 8,100); 1 H NMR (CD 3 OD) δ 3.71 (t, J=7.2 Hz, 2 H, CH 2 N), 3.79 (d, J=9.1 Hz, 2 H, CH 2 P), 3.92 (m, 5 H, C 2 OCH 3 +CH 2 OP ), 4.10, (t, J=7.2 Hz, 2 H, CH 2 O), 4.18 (s, 3 H, C 3 OCH 3 ), 5.50 (t, J=6.8 Hz, 1H, ═CH), 8.76 (s, 1 H, HC 8 ); MS m/z 457 (M + ). Anal. (C 16 H 20 N 5 O 9 P) C, H, N; calcd (%): 42.02, 4.41, 15.31; found (%): 42.14, 4.50, 15.25. EXAMPLE 8 Enzyme Assays For comparison, compounds 1, 9-[(2-hydroxy-ethoxy)methyl]guanine (acyclovir 2), 9-(β-D-arabinofaranosyl)adenine (ara-A 3), 9-[2-(phosphonomethoxy)ethyl]adenine (PMEA, 4), and 9-[2-(phosphonomethoxy)ethyl]guanine (PMEG 27) were either synthesized or purchased from a commercial source. 9-[(2-Hydroxyethoxy)methyl]adenine 1. Compound 1 was prepared by an standard procedure: see Davari (1989) Doctoral Thesis, School of Veterinary Medicine, Shiraz University, Shiraz, Iran. Mp 198-199° C.; R f (hexanes/EtOAc=1:2) 0.23; UV (EtOH) λ max 259 (ε 14,000); 1 H NMR (CD 3 OD) δ 3.62 (s, 4 H, O(CH 2 ) 2 O), 5.67 (s, 2 H, H 2 C 1′ ), 8.22 (s, 1 H, HC 2 ), 8.27 (s, 1 H, HC 8 ); 13 C NMR (CD 3 OD) δ 61.86 (CH 2 OH), 72.08 (OCH 2 ) 74.31 (C 1′ ), 119.98 (C 5 ), 143.14 (C 8 ), 150.90 (C 4 ), 153.74 (C 2 ), 157.19 (C 6 ); MS m/z 209 (M + ). Lipophilicity and Solubility Tests. Lipophilicity and water solubility were determined by the distribution between 1-octanol and water according to the methods reported by Baker et al. (1978 , J. Med. Chem . 21: 1218-1221). The results show in Table 1. To determine lipophoilicity (Partition Coefficients), a solution of each compound (10 mL) in phosphate buffer (0.10 M) possessing an UV absorbance of 2.2-3.3 at 258-267 nm, for adenine, or at 270-290 nm, for guanine, was shaken with 1-octanol (10 mL) in a separatory funnel for 1.5 h. The layers were separated and their concentrations were determined by an UV spectrophotometer. The partition coefficient was calculated as P=[S] 1-octanol /[S] H 2 O . To determine solubility, each compound (70 mg) was agitated in a 25-mL volumetric flask with phosphate buffer (0.10 M, pH 6.8, 5.0 mL) for 20 h. This solution was filtered from undissolved solid through a sintered glass funnel (4.0-5.5 mesh ASTM) and the concentration of the solution was determined by UV absorbance. The results show that adenine acyclic nucleosides 1 and 14 were observed to exhibit higher lipophilicity as well as water solubility than those exhibited by N 7 -guanine acyclic nucleoside 9, acyclovir 2, and ara-A 3. Phosphonate 20 and butenolide ester derivatives 25, 28, and 29 also exhibited higher lipophilicity and water solubility as compared to PMEA 4 or PMEG 27. On the other hand, even though the solubility of nucleotide analog 24 in water was found to be higher than that of PMEA 4, its lipophilicity was lower than that of 4. Furthermore, butenolide ester derivative 25 showed higher lipophilicity as well as water solubility with respect to the parent nucleotide 5′-monophosphate 24. Kinetic Studies of Competitive Inhibition ofAdenosine Deaminase by Acyclic Nucleosides and Nucleotides. The rates of deamination of N 9 -alkylated adenine 1, ara-A, 3, PMEA 4, N 7 -alkylated adenine 14, phosphonate 20, dinucleotide 5′-monophosphate 24, its butenolide ester derivative 25, and PMEA-containing butenolide 28 in the presence of calf mucosal adenosine deaminase (ADA, EC 3.5.4.4) in buffer solutions were determined. See, e.g., Ogilvie et al. (1984) Can. J. Chem . 62: 241-252 and references cited therein. Additionally, the inhibition studies on these compounds were carried out based on the Kaplan method (Table 2). See, e.g., Moosavi-Movahedi et al. (1993) Int. J. Biol. Macromol . 15: 125-129 and references cited therein. The results (Table 2) showed that both the N 9 - and N 7 -acyclic nucleosides 1 and 14 functioned as ADA substrates. The V max of 14 was less than that of 1 by a factor of 4. Compounds 1, 14, and 24 showed competitive inhibition of ADA when ara-A was used as a substrate. However, N 7 -isomer 14 was found to be more efficient than the N 9 -isomer 1 and nucleotide analog 24 as an inhibitor of ADA. PMEA 4, acyclic nucleoside phosphonate 20, and nucleotide-containing butenolides 25 and 28 were neither a good substrate nor an inhibitor of the enzyme. Nucleotide analog 24, however, was a substrate for ADA, but its V max was about 95% less than that of ara-A 3. The slow rate of deamination of compound 24 by ADA may reflect the lack of substrate activity of the acyclic nucleoside phosphonate moiety therein in the active site of the enzyme. Comparison of Phosphorylation of Nucleosides by HSV and Vero Cell Thymidine Kinases. Phosphorylation of nucleosides with HSV or Vero cell thymidine kinase was studied as described previously. See, e.g., Keller et al. (1981) Biochem. Pharmacol . 30: 3071-3077 and references cited therein. Results are summarized in Table 3. The rate of phosphorylation of N 9 -adenine acyclic nucleoside 1, acyclovir 2, ara-A 3, N 7 -guanine acyclic nucleoside 9, and N 7 -adenine acyclic nucleoside 14 as well as a mixture of 2 and 14 (w/w=1:1), a mixture of 1 and 2 (w/w=1: 1), and a mixture of 2 and 3 (w/w=1:1) in the presence of HSV or Vero cell thymidine kinase was determined and the results were compared with those of thymidine. It was observed that N 7 -guanine nucleoside 9, unlike acyclovir, can be phosphorylated by both the HSV thymidine kinase and the host cell kinase. On the other hand, these enzymes were found not to phosphorylate adenine nucleosides 1, 14, and ara-A 3; yet the N 7 -adenine nucleoside 14 induced a 2-fold increase on the rate of phosphorylation of acyclovir 2 (see Table 3). N 9 -adenine nucleosides 1 and 3, however, did not induce any increase on the rate of phosphorylation of 2. It is believed that the N 7 -adenine nucleoside 14 binds (Km=30 mM) to a specific receptor at the active site of the enzyme to exhibit the observed activatory property toward the HSV thymidine kinase. Enzymatic Conversion of Acyclic Nucleoside Phosphonates and Nucleotide 5′-Monophosphate Analog to Their Antivirally Active Diphosphates (Triphosphate Equivalents). Substrate affinities of acyclic nucleoside phosphonates 4, 20, 21, and dinucleotide analog 24 for PRPP synthetase as well as their inhibitory effect against the enzyme were evaluated according to the established procedures. See, e.g., Balzarini et al. (1991) J. Biol. Chem . 266: 8686-8689 and references cited therein. Results are illustrated in Table 4. The assays were then terminated after 4 h by the addition of MeOH. HPLC on an anion-exchange Partisphere column was used to analyze the formation of ATP 21pp, PMEApp, 4pp, 20pp, and 24pp. The substrate affinity of N 7 -acyclic nucleoside phosphonate 20 (Km=4.8 mM) for the enzyme was found to be 4-times less than that of PMEA 4 (Km=1.5 mM). Similarly, the V max for conversion of 20 to 20pp is at least 7-times lower than that for the conversion of 4 to 4pp. On the other hand, nucleotide analog 24 was found to be a better substrate than PMEA 4 for PRPP synthetase (Km=0.57 mM). It can also be phosphorylated (V max =12) by the enzyme at a rate similar to that of AMP 21 (V max =14). It has been hypothesized (Balzarini supra) that the primary amino group at the C-6 position of purines is essential for hydrogen bonding with the enzyme. The NH 2 group in 20 is sterically hindered by an adjacent side chain at the N-7 position. As such, it cannot as effectively interact with the active site of the enzyme when compared to PMEA 4, natural substrate AMP 21, or nucleotide analog 24. Activity ofSnake Venom and Spleen Phosphodiesterases Against Nucleotide Analogs. Snake venom phosphodiesterase (200 units) was dissolved in tris(hydroxymethyl)aminomethane buffer (1.0 mL), which was adjusted to pH 9.2 with 0.1 N HCl. The enzyme solution (0.10 mL) was added to the nucleotide 24 or 25 (0.70 mg) and the mixture was incubated at 37° C. for 8 h. The solution was then applied to Whatman 3-mm paper as a band, which was developed with a mixture of i-PrOH, conc. NH 4 OH, and H 2 O (9:1:2). Degradation products or unreacted starting materials were separated. Dinucleotide 5′- monophosphate 24 gave 9-(β-D-furanosyl)adenine-3′-[[1-(adenin-7-yl-ethoxy)methyl]phosphonate] in about 80% yield. On the other hand, dinucleotide-containing butenolide 25 was completely resistant to the enzyme. Spleen phosphodiesterase (20 units) was dissolved in sodium pyrophosphate buffer (0.01 M, 1.0 mL), which was adjusted to pH 6.5 with phosphoric acid. Nucleotide 24, 25, or 9-(β-D-furanosyl)adenine-3′-[[1-(adenin-7-yl-ethoxy)methyl]phosphonate] (0.70 mg) was dissolved in ammonium acetate buffer (0.05 M, 0.20 mL), which was adjusted to pH 6.5 with acetic acid. An aliquot of the enzyme solution (0.1 mL) was added to the nucleotide solution and the mixture was incubated at 37° C. for 8 h. The solution was then applied to Whatman 3-mm paper as a band and developed with a mixture of i-PrOH, conc. NH 4 OH, and H 2 O (9:1:2). Bands containing products were cut out, which were eluted with H 2 O and the resultant mixture was freeze-dried. The isolated products were characterized by comparison with authentic samples. Dinucleotide 24 afforded phosphonate 20 and adenosine 5′-monophosphate 21 in about 60% yield. 9-(β-D-Furanosyl)adenine-3′-[[1-(adenin-7-yl-ethoxy)methyl]-phosphonate] gave adenosine (40% yield) and 20 (60% yield). Dinucleotide analog 25, having a butenolide ester unit, was found to be stable to the enzyme.Compound 24, possessing both the skeletons of phosphonate 20 and 5′-adenosine monophosphate 21, was dephosphorylated at the 5′-position by snake venom in 80% yield after 8 h. The resultant 9-(b-D-furanosyl)adenine-3′-[[1-(adenin-7-yl-ethoxy)methyl]phosphonate] was hydrolyzed further in the presence of spleen phosphodiesterase to afford adenosine and phosphonate 20 in about 40-60% yield after 8 h. See, Hakimelahi et al. (1995) J. Med. Chem . 38: 4648-4659 and references cited therein. Spleen phosphodiesterase also degraded compound 24 to give 5′-adenosine monophosphate 21 and phosphonate 20 in an overall yield of 60% after 8 h. These results indicate that the phosphodiesterases recognized dinucleotide analog 24 as a normal substrate. In addition, it was found that nucleotide-containing butenolide 25 was completely resistant to snake venom and spleen enzymes. EXAMPLE 9 In Vitro Assays Anti-DNA Virus Activity in Vitro. The newly synthesized compounds were tested for inhibition of cytopathogenicity of the herpes simplex type 1 virus (HSV-1), herpes simplex type 2 virus (HSV-2), thymidine kinase-positive (TK + ) and thymidine kinase-deficient (TK − ) strains of varicella-zoster virus (VZV), and human cytomegalovirus in Vero cell culture up to a level as high as 128 mg/mL. Compounds tested include N 7 -alkylated purines 7, 9, 12, 14, 20, as well as a mixture of 1 and 14 (w/w=1:1), 2 and 14 (w/w=1:1), 3 and 14 (w/w=1:1) in addition to N 9 -substituted adenine 1, acyclovir 2, ara-A 3, PMEA 4, PMEG 27, nucleotide analog 24, and nucleotide-containing butenolides 25, 28, and 29. Toxicity of these compounds was evaluated by their ability to cause morphological changes in HeLa and Vero cells at different concentrations. The minimum inhibitory concentrations (IC 50 ) were measured by use of the linear regression method. See, e.g., Armitage (1983) “Statistical Methods in Medical Research,” Blackwell Scientific Publications, Oxford, U.K.; and Hakimelahi et al. (1990) J. Sci. Iran 1: 186-191. The results are summarized in Table 5. The pronounced anti-DNA virus activity of 14 with respect to the corresponding chloro derivative 12 showed that the presence of a hydroxyl group is essential for antiviral activity. Results from the biological tests also indicated that adenine nucleoside 14 was not effectively deaminated by ADA; yet it inhibited the deactivating property of the enzyme and led to the observed increase in the antiviral activity of 1 and ara-A 3. Furthermore, use of adenine acyclic nucleoside 14 resulted in a 2-fold increment in the rate of phosphorylation of acyclovir 2 by HSV-thymidine kinase. As such, a combination of 14 and 2 exhibited profound antiviral activity. N 7 -Adenine nucleoside 14 was found to be less toxic than the corresponding N 9 -isomer 1. On the other hand, both HSV and cellular thymidine kinases can phosphorylate N 7 -guanine nucleoside 9. As a result, this nucleoside exhibited more toxicity than acyclovir 2. The rate of phosphorylation of N 7 -acyclic nucleoside phosphonate 20 to its antivirally active anabolite 20pp by PRPP synthetase is 7-times less than that of the PMEA 4. Thus in comparison to compound 4, compound 20 exhibited less activity against DNA-viruses. On the other hand, nucleotide 5′-monophosphate analog 24, possessing a natural AMP moiety, was converted to its diphosphate (triphosphate form) 24pp at a rate comparable to that of AMP 21, which is about 120-times faster than the rate of conversion of PMEA 4 to PMEApp 4pp. Consequently, nucleotide 5′-monophosphate 24 exhibited higher anti-DNA virus activity than PMEA 4. The ability of a drug to penetrate a membrane and exhibit biological activity can be correlated to its lipophilicity. See, e.g., Hakimelahi et al. (1995). J. Med. Chem . 38: 4648-4659 and references cited therein. Consequently, compounds 25, 28, and 29 possessed butenolide ester functionalities as lipophilic prodrugs. These compounds displayed superior antiviral activity relative to their respective parent compounds nucleotide 5′-monophosphate 24, PMEA 4, and PMEG 27. In addition, spleen phosphodiesterase can recognize and at least partly hydrolyze nucleotide analog 24 to the biologically less active phosphonate 20 inside the infected cells; whereas its butenolide ester derivative 25 was found to be stable toward phosphodiesterases. Thus in comparison with nucleotide analog 24, its prodrug 25 possesses superior bioavailability and greater stability both in vitro and in vivo. Anti-Retrovirus Activity in Vitro. The methods for measuring viruses-induced cytopathogenicity in MT4 cells or CEM cells, as well as the toxicity of the tested compounds towards MT4 and CEM cells have been described previously. See, e.g., Averett (1989) J. Virol. Methods 23: 263-276. Results are summarized in Table 6. Compounds 4, 20, 24, 25, 27, 28, and 29 were tested for inhibition of cytopathogenicity against the human immunodeficiency viruses HIV-1 (III-B) and HIV-2 (LAV-2) in MT4 cells. These compounds were also screened for their antiviral activity against moloney murine sarcoma virus (MSV) in CEM cells in a cell-protection assay. 30 Toxicity of these compounds was evaluated by their ability to cause morphological changes in MT4 or CEM cells at different concentrations. The minimum inhibitory concentrations (IC 50 ) were measured by the use of the linear regression method (Table 6). In comparison to the rate of phosphorylation of PMEA 4 to PMEApp 4pp by PRPP synthetase, the conversion of dinucleotide 24 to its anabolically active form 24pp is 120 times faster; yet PMEA 4 exhibited higher activity than 24 as well as the butenolide ester derivative 25 against retroviruses. Thus, the HIV and MSV reverse transcriptases may have higher affinity for PMEA 4 than dinucleotide analog 24. PMEA 4 was also found to be more active than its N 7 -isomer 20 against retroviruses. On the other hand, butenolide ester derivatives 28 and 29 displayed superior antiviral activity relative to their respective parent molecules 4 and 27. Thus in comparison to PMEA 4 and PMEG 27, their respective lipophilic prodrugs 28 and 29 possess superior bioavailability and greater anti-retrovirus activity. As shown in Scheme 6, we believe that the oxygen of the methoxy group at the C-2 position of the butenolide moiety is responsible for the ease of conversion of these novel prodrugs 28 and 29 to their corresponding potential drugs PMEA 4 and PMEG 27 inside the infected cells. EXAMPLE 10 In Vivo Assays Anti-HSV-1 Activity in Vivo and Determination of LD 50 for N 7 -Adenine Acyclic Nucleoside 14, Nucleotide-Containing Butenolide 25, and PMEA-Containing Butenolide 28 in Mice. Two-weeks-old NMRI mice (15-20 animals/group), weighing ca. 7 g each, were infected i.p. with 4×10 4 units of HSV-1 (KOS). See, e.g., Kim et al. (1991) J. Med. Chem . 1991, 34, 2286-2294. Compounds in Table 7 were administered i.p. once a day for 6 consecutive days, starting 4 h postinfection. Percentage of HSV-1-infected mice without symptoms and those that were alive at day 21 postinfection were observed (see Table 7). Deaths were recorded for 21 days after infection. Acyclovir 2, PMEA 4, N7-guanine acyclic nucleoside 9, N 7 -adenine acyclic nucleoside 14, N 7 -acyclic nucleoside phosphonate 20, nucleotide 5′-monophosphate 24, and butenolide ester derivatives 25 and 28 were evaluated for their inhibitory effect on HSV-1-induced mortality in NMR 1 mice (Table 6). Butenolide derivative of PMEA, 28, appeared to be the most potent anti-HSV-1 agent in vivo, followed by nucleotide-containing butenolide 25, nucleoside analog 14, nucleotide analog 24, acyclovir 2, PMEA 4, phosphonate 20, and nucleoside analog 9. Since compound 28 is less active in vitro against HSV-1 when compared to compounds 25 and 14 respectively, the in vitro potency does not directly translate to in vivo potency. These results confirmed previous findings. See, e.g., De Clercq et al. (1986) Nature (London) 323: 464-467; and Pauwels et al. (1988) Agents Chemother . 32: 1025-1030. All compounds were administered intraperitoneally (i.p., 100-250 mg/kg/day) for 6 consecutive days. Compounds 2, 4, 14, 24, and 25 gave full protection against HSV-induced mortality at the 150 mg/kg dose level. The same level of protection was provided by compound 28 at a dose of 100 mg/kg. Survival times of all treated groups were found to be significantly different from the placebo treated control group (see Table 7). The potent anti-HSV-I activity exhibited by compounds 2, 4, 14, 24, 25, and 28 clearly demonstrated that they are taken up effectively by cells to exert in vivo activity. None of the compounds were toxic to the mice at the highest dose tested. The LD 50 values of the most active compounds 14, 25, and 28 in mice were also determined. As such, N 7 -acyclic nucleoside 14 and butenolide ester derivatives 25 and 28 were administered at different doses intraperitoneally. They did not show any toxicity up to a concentration level as high as 400 mg/kg. All mice were controlled in good conditions after six months of administration. Nevertheless, LD 50 (i.p.) values of 950 mg/kg, 675 mg/kg, and 710 mg/kg were determined for 14, 25, and 28, respectively. Moreover, no discernible abnormality was observed in the histological appearance of the viscera of either the control or tested groups of mice that received the drugs i.p. (250 mg/kg/day) for 10 days. Furthermore, there were no physiological changes in their cardiovascular or central nerve systems. Inhibitory Effects of PMEA 4 and Its Butenolide Ester Derivative 28 on MSV-Induced Tumor Formation in Vivo. The inhibitory effects of the compounds 4 and 28 on the initiation of MSV-induced tumor formation and survival of MSV-induced mice (10-15 animals/group) were evaluated as previously described. See, e.g., Balzarini et al. (1993) Agents Chemother 37: 332-338 and references cited therein. Results are summarized in Table 8. Compounds 4 and 28 were evaluated for their inhibitory effect on MSV-induced tumor formation in NMR 1 mice (Table 8). The compounds were administered intraperitoneally (50 mg/kg/day) for two consecutive days. Prodrug 28 exhibited much higher anti-MSV activity than PMEA 4 in vivo. At a dose of 10 mg/kg/day, compound 28 prevented tumor formation in 60% of the MSV-infected mice whereas with compound 4 at the same dosage level, only 19% prevention was observed. In surviving animals treated with 28, about 2 g weight loss was observed. In the case of PMEA-treated mice, the weight loss of the surviving animals was at least two times more. TABLE 1 Solubility and Lipophilicity of Nucleoside and Nucleotide Analogs solubility in log P compound water (mg/mL) (1-octanol/water) a  1 2.56 0.98 acyclovir 2 0.40 −0.48 ara A 3 0.48 −0.50 PMEA 4 1.97 0.69  9 0.32 −0.60 14 4.71 1.24 20 2.89 0.95 24 3.08 0.12 25 5.16 1.37 PMEG 27 0.36 0.14 28 7.42 2.09 29 2.05 0.79 a Parition coefficients were calculated as P = [substrate] 1-octanol /[substrate] H 2 O . TABLE 2 Substrate Activity and Inhibitory Property Against ADA a Substrate K m (□M) rel. V max K i (μM)  1 138.6 1.48 × 10 −2 140.8 ara A 3 42.8 1.0 — PMEA 4 427.0 — >800 14 198.5 1.50 × 10 −6 8.3 20 >800 — >800 24 164.5 9.78 × 10 −2 99.7 25 635.7 — >800 28 >800 — >800 a The reaction velocity, V, in μmol/min/mg of the enzyme was determined, and a plot of 1/[S] ([S] = substrate concentration) vs. 1/V was made. Thus, by the method of Lineweaver and Burk, the Michaelis (K m ) and maximum velocity constants (V max ) were determined. For inhibition studies, in addition to ara-A 3, substrate solutions contained 1, 4, 14, 20, 24, 25 or 28 and then K i was measured for each substrate. TABLE 3 Phosphorylation of Various Nucleosides and Thymidine with HSV or Vero Cell Thymidine Kinases a HSV thymidine kinase Vero cell thymidine kinase Substrate Km (□M) rel. Vmax Km μM) rel. Vmax  1 2.0 × 10 4 <3.0 >3.0 × 10 4 <3.0 acyclovir 2 1.5 39.2  2.2 × 10 4 <3.0 ara A 3 1.5 × 10 4 <3.0  1.3 × 10 4 <3.0  9 12.8 28.0 18.3 12.0 14 30.0 <3.0 >3.0 × 10 4 <3.0  2 + 14 (1:1 w/w) — 80.9 — <3.0 thymidine 1.0 100.0 1.0 1.0 × 10 2 a Apparent Km values were determined for HSV and Vero cell thymidine kinases from reactions containing 50 mM of Tris-HCl (pH 7.5), 2.0 mM of ATP, 2.0 mM of MgCl 2 , 1.0 mg/mL of BSA, 1.5 μM of [C 14 ]-thymidine, and 198 units of enzymes/mL. The results were compared with those of thymidine. All reactions were performed at 37° C. The radiochemical nucleoside kinase coupled assay was used in the determination of the relative substrate velocities. TABLE 4 Kinetics of the PRPP Synthetase Reaction with Acyclic Nucleoside Phosphonates 4 and 20, Nucleotide Analog 24, and AMP (21) a substrate K m (mM) V max (μmol/unit/h) K i b (mM) PMEA 4 1.51 0.096 3.02 20 4.83 0.013 16.74 AMP 21 0.24 14.270 — 24 0.57 11.560 0.79 a The PRPP synthetase reaction mixture contained 10.0 mM potassium phosphate buffer (pH 8.0), 5.0 mM MgCl 2 , 2.5 mM PRPP, an appropriate amount of AMP (21) or test compounds, and 0.04 unit of PRPP synthetase. The formation of ATP, diphosphates of PMEA 4, 20, or 24 was analysed by HPLC according to the method of Balzarini and De Clercq. b Inhibition of AMP phosphorylation by the enzyme was measured in the presence of different substrate (AMP) concentrations and appropriate concentrations of 4, 20, and 24. The reaction mixture was incubated at 37° C. for 15 min with 0.002 unit of PRPP synthetase. The formation of ATP was followed by HPLC using a Partisphere anion exchange column. TABLE 5 Anti-DNA Virus and Anticellular Activities of Nucleoside and Nucleotide Analogs in Tissue Culture. IC 50 a (μg/mL) HSV-1 HSV-2 TK + VZV TK − VZV HCMV compound (KOS) (G) (YS) (YS/R) (AD-169) HeLa cell b Vero cell b  1 4.5 8.4 7.0 9.0 4.8 196 209 acyclovir 2 0.46 1.2 9.7 28 32 259 269 ara A 3 8.8 >128 >128 >128 27 73 89 PMEA 4 11 7.4 6.1 6.0 15 146 175  7 >128 >128 >128 >128 >128 157 145  9 1.1 2.8 13 38 16 114 99 12 >128 >128 >128 >128 >128 169 165 14 0.58 0.97 1.1 0.97 1.0 472 496 20 20 13 8.5 8.0 21 236 246 24 7.0 4.1 3.8 4.2 11 244 236 25 3.0 1.9 1.7 2.0 4.1 326 350 PMEG 27 4.6 6.3 0.040 0.060 0.35 9.1 7.0 28 3.7 3.0 2.0 2.4 5.4 139 170 29 0.42 0.28 0.0070 0.0080 0.020 8.8 6.6  1 + 14 c 0.060 0.10 0.17 0.23 0.18 372 388  2 + 14 c 0.010 0.020 0.49 0.71 0.06 420 410  3 + 14 c 0.10 1.1 0.98 1.7 0.46 280 293 a Compound concentration required to inhibit virus-induced cytopathogenicity by 50% was determined according to an established procedure. IC 50 values represent the mean of duplicate determinations. b Concentration of the compound required to cause microscopically visible change or disruption in about 50% of the cell sheet. c For 1:1 (W/W) mixtures, IC 50 in μg/mL refers to known antiviral drugs. TABLE 6 Inhibitory Effects of Nucleotide Analogs on the Cytopathogenicity of HIV-1 and HIV-2 in MT4 Cells, as well as on the Cytopathogenicity of MSV in CEM Cells and Cellular Toxicity. IC 50 a (μg/mL) MT4 CEM compound HIV-1 (IIIB) HIV-2 (LAV-2) MSV cell b cell b PMEA 4 4.1 3.8 2.0 274 285 20 7.8 9.1 27 >300 >300 24 5.9 6.4 17 298 >300 25 4.9 4.2 13 >300 >300 PMEG 27 16 18 0.19 16 12 28 1.4 1.0 0.93 265 280 29 6.0 7.1 0.020 14 13 a Inhibitory concentrations (IC 50 ) were determined by use of an established procedure and represent the average of duplicate determinations. b Concentration of the compound required to reduce the number of viable uninfected cells by 50%. TABLE 7 Antiviral Effects of compounds 2, 4, 9, 14, 20, 24, 25, and 28 Against HSV-1-Induced Mortality in NMRI Mice Upon Intraperitoneal Administration a dose No. mean mean (mg/kg/ of day of symptom day of animal compound day) mice initiation (%) b death (%) c acyclovir 2 250 20 >21 (100%) >21 (100%) 150 20 19.1 ± 1.3 (86%) >21 (100%) 100 15 15.6 ± 1.6 (65%) 18.9 ± 2.1 (80%) PMEA 4 250 20 >21 (100%) >21 (100%) 150 20 18.5 ± 1.9 (80%) >21 (100%) 100 15 14.7 ± 1.4 (56%) 17.0 ± 1.1 (77%)  9 250 20 15.1 ± 1.2 (67%) 18.5 ± 1.0 (90%) 150 20 13.0 ± 1.1 (57%) 16.0 ± 1.5 (78%) 100 15 10.9 ± 0.6 (48%) 14.1 ± 0.8 (66%) 14 250 20 >21 (100%) >21 (100%) 150 20 19.9 ± 1.3 (88%) >21 (100%) 100 15 17.5 ± 2.1 (80%) 19.8 ± 1.7 (95%) 20 250 20 16.8 ± 2.4 (75%) 20.0 ± 1.7 (95%) 150 20 15.0 ± 0.9 (62%) 17.8 ± 1.5 (80%) 100 15 12.8 ± 1.1 (51%) 15.4 ± 1.3 (70%) 24 250 20 >21 (100%) >21 (100%) 150 20 19.3 ± 1.6 (85%) >21 (100%) 100 15 16.0 ± 1.7 (70%) 19.5 ± 1.2 (90%) 25 250 20 >21 (100%) >21 (100%) 150 20 >21 (100%) >21 (100%) 100 15 19.6 ± 1.5 (94%) >21 (100%) 28 250 20 >21 (100%) >21 (100%) 150 20 >21 (100%) >21 (100%) 100 15 >21 (100%) >21 (100%) Saline 0 20 3.38 ± 0.7 (0%) 9.4 ± 0.6 (0%) a Mice were inoculated intraperitoneally with HSV-1 (KOS). Treatment was initiated 4 h postinfection and continued for 6 consecutive days. Experiments were terminated at day 21. b Values in parantheses represent percentage of HSV-1-infected mice without symptoms at day 21 postinfection. c Values in parantheses represent percentage of HSV-1-infected mice that were alive at day 21 postinfection. TABLE 8 Inhibitory Effects of Acyclic Nucleoside Phosphonates 4 and Its Prodrug 28 on MSV-Induced Tumor Formation and Associated Death in NMRI Mice Upon Intraperitoneal Administration a dose No. (mg/kg of mean day of tumor mean day of animal compound day) mice initiation (%) b death (%) c PMEA 4 50 15 12.5 ± 1.6  (84%) 17.8 ± 2.0 (96%) 20 15 12.0 ± 1.3  (58%) 15.9 ± 1.7 (76%) 10 10 9.6 ± 1.5  (19%) 13.0 ± 1.1 (42%) 28 50 15 18.4 ± 1.4  (95%) >21 (100%) 20 15 17.9 ± 1.7  (80%) >21 (100%) 10 10 14.0 ± 2.1  (60%) 18.9 ± 1.8 (97%) control 0 30 3.82 ± 0.95 (0%) 7.8 ± 1.3 (0%) untreated control d 0 40 >21 (100%) >21 (100%) a All mice received two injections within two days. b Values in parantheses represent percentage of MSV-infected mice without tumors at day 21 postinfection. c Values in parantheses represent percentage of MSV-infected mice that were alive at day 21 postinfection. d Untreated control group was neither treated with MSV nor with the drugs. Other Embodiments All of the features disclosed in this specification may be combined in any combination. Each feature disclosed in this specification may be replaced by an alternative feature serving the same, equivalent, or similar purpose. Thus, unless expressly stated otherwise, each feature disclosed is only an example of a generic series of equivalent or similar features. From the above description, one skilled in the art can easily ascertain the essential characteristics of the present invention, and without departing from the spirit and scope thereof, can make various changes and modifications of the invention to adapt it to various usages and conditions. Thus, other embodiments are also within the claims.

This invention relates to purine compounds of formula (I): R 1 is NH 2 or OH; R 2 is H or NH 2 ; R 3 is H or alkyl; each of m and n, independently, is 1, 2, 3, or 4; X is O, S, or NH; and Y is H, halogen, OR a , P(O)(OR a ) 2 , or P(O)(OR a )(OR b ), in which R a is H, alkyl, aryl, heteroaryl, cyclyl, heterocyclyl, and R b is wherein A is adenine, guanine, cytosine, uracil, or thymine; R c is H or OH; R d is H or alkyl; R e is H, alkyl, or 5-ethylidene-(3,4-dialkoxyl)-furan-2-one; provided that if R 1 is NH 2 , R 2 is H; and if R 1 is OH, R 2 is NH 2 .

PRIOR RELATED APPLICATIONS This application is a national stage application of PCT Patent Application PCT/CU2004/000002, filed Feb. 19, 2004, which claims priority to Cuban patent applications CU 2003-0039, filed Feb. 20, 2003, and CU 2003-0084, filed Apr. 17, 2003. CROSS REFERENCE This application is a national stage of PCT/CU2004/000002 filed Feb. 19, 2004 and is based upon Cuban Patent Applications No. 2003-0039, filed Feb. 20, 2003 and No. 2003-0084 filed Apr. 17, 2003 under the International Convention. FIELD OF THE INVENTION The field of invention is that of biotechnology, in particular, the obtainment of Vibrio cholerae live attenuated vaccine strains, more specifically, the introduction of defined mutations to prevent or limit the possibility of reacquisition and (or) the later dissemination of CTXΦ phage encoded genes by those live vaccine strains and a method to preserve them to be used as vaccines. BACKGROUND OF THE INVENTION First, Definitions: During the description of the invention will be used a terminology whose meaning is listed bellow. By CTXΦ virus is meant the particle of protein-coated DNA produced by certain V. cholerae strains, which is capable of transducing its DNA, comprising cholera toxin genes, to other vibrios. By cholera toxin (CT) is meant the protein responsible for the clinical symptoms of cholera when produced by the bacteria. By CTXΦ-encoded toxin genes are meant, in addition to CT genes, zot and ace genes that encode for the “zonula occludens toxin” and for the accessory cholera enterotoxin, respectively. The activity of ZOT is responsible for the destruction of the tight junctions between basolateral membranes of the epithelial cells and ACE protein has an activity accessory to that of the cholera toxin. The term well tolerated vaccine or well tolerated strain refers to such strain lacking the residual reactogenicity that characterize most of the of non-toxigenic strains of V. cholerae . In practical terms, it means that it is a strain safely enough to be used in communities without or with limited access to healthcare institutions without risks for the life of the vaccinees. It should be expected a rate of diarrhea in 8% or less of the vaccinees and the diarrhea is characterized in that it does not exceed 600 ml (grs), only 1% of the vaccinees or less could suffer from headache, which should be minor and of short duration (less than 6 h), and finally that it prompts vomits in less than 0.1% of the vaccinees, those vomits characterized for being a single episode of 500 ml or less. By hemagglutinin protease (HA/P) is meant the protein secreted by V. cholerae manifesting dual function, being one of them the ability to agglutinate the erythrocytes of certain species and the other the property to degrade or to process proteins such as mucine and the cholera toxin. By celA is meant the nucleotide sequence coding for the synthesis of the endoglucanase A. This protein naturally occurs in Clostridium thermocellum strains and has a β (1-4) glucan-glucano hidrolase activity able to degrade cellulose and its derivatives. The term MSHA is referred to the structural fimbria of the surface of V. cholerae with capacity to agglutinate erythrocytes of different species and that is inhibited by mannose. By reversion to virulence mediated by VGJΦ is meant the event in which a previously attenuated strain obtained by the suppression of CTXΦ genes reacquire all the genes of this phage through a mechanism completely dependent and mediated by VGJΦ and the interaction with its receptor, MSHA. The possibility of disseminating the CTXΦ phage in a process mediated by VGJΦ is that in which the filamentous phage VGJΦ form a stable hybrid structure (HybPΦ) through genetic recombination with the DNA of CTXΦ and disseminate its genome with active genes toward other strains of V. cholerae , which could be environmental non pathogenic strains, vaccine strains or other from different species. Second, information of the previous art: Clinical cholera is an acute diarrheal disease that result from an oral infection with the bacterium V. cholerae . After more than 100 years of research in cholera there remains the need for an effective and safe vaccine against the illness. Since 1817 man has witnessed seven pandemics of cholera, the former six were caused by strains of the Classical biotype and the current seventh pandemic is characterized by the prevalence of strains belonging to El Tor biotype. Recently, beginning in January of 1991, this pandemic extended to South America, and caused more than 25 000 cases of cholera and over 2 000 deaths in Peru, Ecuador and Chile. By November 1992, a new serogroup of V. cholerae emerged in India and Bangladesh, the 0139, showing a great epidemic potential and generating great concern through the developing world. These recent experiences reinforce the need for effective cholera vaccines against the disease caused by V. cholerae of serogroups O1 (biotype El Tor) and O139. Because convalescence to cholera is followed by an state of immunity lasting at least three years, much efforts in Vibrio cholerae vaccinology have been made to produce live attenuated cholera vaccines, that closely mimics the disease in its immunization properties after oral administration, but do not result reactogenic to the individuals ingesting them (diarrhea, vomiting, fever). Vaccines of this type involve deletion mutations of all toxin genes encoded by CTXΦ. For example, the suppression of the cholera toxin and other toxins genes encoded in the prophage CTXΦ is a compulsory genetic manipulation during the construction of a live vaccine candidate (see inventions of James B. Kaper, WO 91/18979 and John Mekalanos WO 9518633 of the years 1991 and 1995, respectively). This kind of mutants have been proposed as one dose oral vaccines, and although substantially attenuated and able to generate a solid immune responses (Kaper J. B. and Levine M. Patentes U.S. Pat. Nos. 06,472,276 and 581,406). However, the main obstacle for the widespread use of those mutants has been the high level of adverse reactions they produce in vaccinees (Levine and cols., Infect. and Immun. Vol 56, No1, 1988). Therefore, achieving enough degree of attenuation is the main problem to solve during the obtainment of live effective vaccines against cholera. There are at least three live vaccine candidates, which have shown acceptable levels of safety, i.e., enough degree of attenuation and strong immunogenic potential. They are V. cholerae CVD103HgR (Classical Biotype, serotype Inaba) (Richie E. and cols, Vaccine 18, (2000): 2399-2410.), V. cholerae Perú-15 (Biotype El Tor, serotype Inaba) (Cohen M., and cols. (2002) Infection and Immunity, Vol 70, Not. 4, pag 1965-1970) and V. cholerae 638 (Biotype El tor, serotype Ogawa) (Benítez J. A. and cols, (1999), Infection and Immunity. February; 67(2):539-45). Strain CVD103HgR is the active antigenic component of a live oral vaccine against cholera licensed in several countries of the world, the strains Perú-15 and 638 are other two live vaccine candidates to be evaluated in field trials in a near future. However, there is a second problem of importance to solve in those live attenuated vaccine candidates; one is the environmental safety, specially related with the possible reacquisition and dissemination of the cholera toxin genes by existent mechanisms of horizontal transfer of genetic information among bacteria. In accordance with this, the attenuated vaccine strains of V. cholerae , could potentially reacquire virulence genes out of the controlled conditions of the laboratory, in an infection event with CTXΦ phage (Waldor M. K. and J. J. Mekalanos, Science 272:1910-1914) coming from other vibrios and later on contribute to their dissemination. This process could become relevant during vaccination campaigns where people ingest thousands of millions of attenuated bacteria and keep shedding similar quantities in their stools during at least 5 days. Once in the environment, bacteria have the possibility of acquiring genetic material from other bacteria of the same or different species of the ecosystem. For these reasons, at present it is desirable to obtain vaccine candidates with certain characteristics that prevent or limit the acquisition and dissemination of CTXΦ, and especially of the genes coding for the cholera enterotoxin. As a consequence, this is the field of the present invention. Bacterial viruses, known as bacteriophages, have an extraordinary potential for gene transfer between bacteria of the same or different species. That is the case of CTXΦ phage (Waldor M. K. and J. J. Mekalanos, 1996, Science 272:1910-1914,) in V. cholerae . CTXΦ the genes of carries the genes that encode cholera toxin in V. cholerae and enters to the bacteria through interaction with a type IV pili, termed TCP, from toxin co-regulated pilus. TCP is exposed on the external surface of the vibrios. In accordance with published results, under optimal laboratory conditions the CTXΦ phage reaches titers of 10 6 particles or less by ml of culture in the saturation phase; this allows classifying it as a moderately prolific bacteriophage. Equally the expression of the TCP receptor of this phage has restrictive conditions for its production. In spite of these limitations, the existence of this couple bacteriophage-receptor, limits in some way the best acceptance of live cholera vaccines, that is why depriving the bacteria from the portal of entrance to this phage is a desirable modification. There are two theoretical ways of preventing the entrance of CTXΦ into V. cholerae, 1) suppressing the expression of TCP or 2) removing the TCP sites involved in phage receptor interaction. None of the two forms has been implemented due to the essentiality of TCP for proper colonization of the human intestine and elicitation of a protective immune response. It should be noted that sites involved in the TCP-CTXΦ interaction are also needed for the colonization process. (Taylor R. 2000. Molecular Microbiology, Vol (4), 896-910). Several strategies that counteract the entrance of the virus have been evaluated such as preventing the integration of the phage to the bacterial chromosome and its stable inheritance, consisting in the suppression of the integration site and in the inactivation of recA gene to avoid recombination and integration to other sites of the chromosome. (Kenner and cols. 1995. J. Infect. Dis. 172:1126-1129). Also, it has been recently described that the entry of CTXΦ into V. cholerae depends on the genes TolQRA, however this mutation produces sensitive phenotypes not undesired in vaccine candidates of cholera and it has not been implemented. (Heilpern and Waldor. 2000. J. Bact. 182:1739). Further methods that prevent the entrance of phages carryings essential virulence determinants to cholera vaccine strains or other vaccine strains have not been described. SUMMARY OF THE PRESENT INVENTION The main subject of the present invention is related with the phage VGJΦ and its capacity to transfer the genes coding for the cholera toxin, using the Mannose Sensitive Hemagglutinin (MSHA) fimbria as receptor. Specifically, it consists in protecting the live attenuated vaccine strains from the infection with VGJΦ by introducing suppression mutations or modifications that prevent the correct functioning of this fimbria. In the previous knowledge of this fimbria, the following aspects can be summarized. The gene product of mshA was originally described to be the major subunit of a fimbrial appendage in the surface in V. cholerae that had the capacity to agglutinate erythrocytes of different species, this capacity being inhibited by mannose (Jonson G. and cols (1991). Microbial Pathogenesis 11:433-441). As such, the MSHA was considered a virulence factor of the bacteria (Jonson G. and cols (1994). Molecular Microbiology 13:109-118). In accordance with the attributed importance, mutants deficient in the expression of the MSHA were obtained to study its possible role in virulence. It was demonstrated that MSHA, contrary to TCP, is not required for colonization of the human small intestine by the El Tor and O139 V. cholerae (Thelin KH and Taylor RK (1996). Infection and Immunity 64:2853-2856). The MSHA has been also described as the receptor of the bacteriophage 493 (Jouravleva E. and cols (1998). Infection and Immunity, Vol 66, Not 6, pag 2535-2539), suggesting that this phage could be involved in the emergence of the O139 vibrios (Jouravleva E. and cols, (1998). Microbiology 144:315-324). Later on it has been described that the fimbria MSHA has a role in biofilm formation on biotic and a-biotic surfaces contributing thus to bacterial survival outside of the laboratory and the host (Chiavelli D. A. and cols, (2001). Appl. Environ Microbiol. July; 67(7):3220-25 and Watnick P. I. and Kolter R. (1999). Mol. Microbiol. November, 34(3):586-95). It is evident from the previous data that several investigations related with the MSHA fimbria have been done, but none of them defines this pili as the receptor of a phage able to transduce in a very efficient way the genes of the cholera toxin and not only these genes but the complete genome of CTXΦ, what could notably contribute to their dissemination. Additionally, although an extensive search has been made no inventions related with this fimbria have been found, either as virulence factor or as a phage receptor mediating dissemination of CTXΦ. On the other hand, it is common practices among those who develop live cholera vaccines to provide them freeze-dried. Thus, these preparations of the live bacteria are ingested after the administration of an antacid solution that regulates the stomach pH and so the bacterial suspension continues toward the intestine without being damaged in the stomach and achieves colonization in the intestine. Elaboration of freeze-dried vaccines improves preservation of strains, facilitates preparation of doses, allows a long-term storage, limits the risks of contamination and makes the commercialization and distribution easier, without the need of a cold chain, generally not available in under developing countries. Although Vibrio cholerae is considered a very sensitive microorganism to the freeze-drying process, some additives are known to enhance strain survival. Thus, for preservation of the vaccine strain CVD103HgR Classical Inaba, the Center for vaccine Development, University of Maryland, United States, the Swiss Institute of Sera and Vaccines, from Berne (ISSVB), developed a formulation, see (Vaccine, 8, 577-580, 1990, S. J. Cryz Jr, M. M. Levine, J. B. Kaper, E. Fürer and B) that mainly contain sugars and amino acids. The formulation is composed of sucrose, amino acids and ascorbic acid, and after the freeze-drying process, lactose and aspartame are added. In a work about preservation by freeze-drying of the wild type strain 569B Classical Inaba, published in Cryo-Letters, 16, 91-101 (1995) for Thin H., T. Moreira, L. Luis, H. García, T. K. Martino and A. Moreno, compared the effect of different additives on the viability and final appearance upon liophilization and after the storage at different temperatures of this V. cholerae strain. It was demonstrated that viability losses were less than 1 logarithmic order after 3 days of storage to 45° C. The invention CU 22 847 claims a liophilization method where the formulations contain a combination of purified proteins or skim milk with addition of polymers and/or glycine, besides bacteriologic peptone or casein hydrolysate and sorbitol, with good results for the viability of Vibrio cholerae strains of different serogroups, biotypes and serotypes. The freeze-dried bacteria keep their viability after being dissolved in a 1,33% sodium bicarbonate buffering solution used to regulate the pH of the stomach. Any vaccine formulation of cholera that it is supposed to be used in under developing countries should have certain requisites such as posses a simple composition, be easy to prepare and manipulate, be easy to dissolve and have good appearance after dissolved. Besides, It would be also desirable not to require low storage temperatures and to tolerate high storage temperatures at least for short periods of time, as well as the incidental presence of oxygen and humidity in the container. Additionally, it is also necessary an adequate selection of the composition of the formulation that allows the preservation of Vibrio cholerae of different serogroups, biotypes and serotypes. Finally, it is also remarkable that a formulation free of bovine derivate ingredients allows us to be in agreement with the international regulatory authorities related to the use of bovine components due to the Bovine Spongiform Encephalopathy Syndrome. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 . Microphotography of VGJΦ phage. Magnification×32 000. VGJΦ phage was purified from the supernatants of infected Vibrio cholerae 569B. FIG. 2 . Diagram of the genome of hybrid phage HybPΦ-Kn, which has a high potentiality for cholera toxin transmission. The att sequences shown are SEQ ID NO: 12-15. FIG. 3 . Scheme of the genetic manipulation used to suppress mshA gene of V. cholerae vaccine candidates and the suicide vector used during the proceeding. FIG. 4 . Suckling Mice survival inoculated with an attenuated strain and its derivative infected with HybPΦ-Kn that revert it to virulence. DESCRIPTION OF THE INVENTION The present invention propose a new generation of live attenuated vaccines to immunize against cholera by modification of their properties, specifically improving their biological safety during colonization of humans and later in the environment, outside the laboratories. The present invention born from the necessity to protect live cholera vaccines from infection with the CTXΦ bacteriophage, which contains the cholera toxin genes, and also to impair the potential dissemination of this phage starting from live cholera vaccine candidates. Specifically it was born from the discovery and characterization of the VGJΦ phage in our laboratory. VGJΦ is a filamentous bacteriophage isolated from V. cholerae O139 but it has infective capacity on V. cholerae O1 of all serotypes and biotypes and also over other strains of V. cholerae O139. The sequence of this phage was not described in the complete genome sequence of V. cholerae , indicating that this phage was not present in the strain N16961 (O1, El Tor Inaba). From a broad list of V. cholerae O1 strains existing in our laboratory, none of them had homologous sequences to VGJΦ, while strains MO45, SG25-1 and MDO12C, of V. cholerae O139 had. The VGJΦ phage infects V. cholerae through the MSHA fimbria. When this phage enters to the bacterium it can replicate or integrate into a specific chromosomal region. This is a very active phage that reaches 10 11 particles ml −1 in the culture supernatants. The most important characteristic in this phage, by virtue of which the following application of invention is issued, is their capacity to carry out a specialized transduction of the CTXΦ phage and consequently of the cholera toxin genes. This process occurs by a site-specific recombination between CTXΦ and VGJΦ genome, followed by the encapsulation and exportation of both genomes into the VGJΦ capsid. This hybrid viral particle was named HybPΦ. A culture of bacteria infected with both, CTXΦ and VGJΦ, produce 10 11 particles ml −1 of VGJΦ and 10 7 -10 8 particles ml −1 of HybPΦ, which is at least 100 times higher than the titers obtained with CTXΦ alone. It is also important to understand, to the purpose of this application that the CTXΦ phage receptor is TCP, which require special conditions for its expression, while the VGJΦ receptor is MSHA fimbria, an antigen that is expressed abundantly in all culture conditions studied and that is also produced in the environment. Furthermore, other vibrios produce the MSHA what increase the risk of transmission, even to other bacterial species. It is also important to know that once, a new host become infected with HybPΦ, a stable production of particles in the range of 10 7 -10 8 ml −1 takes place in the saturation phase, thus this hybrid phage has a high potential to transmit and disseminate the cholera toxin genes. Another aspect of supreme interest to the purpose of this invention is that cholera toxin genes in HybPΦ are active enough to produce 50 ng ml −1 of toxin during in vitro culture and that the infection of an attenuated strain with HybPΦ revert it back to virulence as assessed by the infant mouse cholera model. In accordance with these data, a primary objective of the present invention is to describe the additional mutations made to live cholera vaccines to prevent them to be infected with either VGJΦ or HybPΦ, as well as the necessity to use live cholera vaccines from which the genome of VGJΦ is absent to avoid the dissemination of CTXΦ mediated by VGJΦ, in the case of reacquisition of CTXΦ. An example of this mutation is a stable spontaneous mutation, conducive to the lack of expression of the MSHA fimbria in the cellular surface. This way, the VGJΦ or its derivative phage, HybPΦ could not infect such vaccines. Another example of this mutation is a suppressive mutation in the structural gene of the major protein subunit of this fimbria (MshA). The use of live cholera vaccine candidates in which the genome of VGJΦ is absent could be achieved simply by searching for hybridization of DNA in different strains to identify which have not homologous fragments to VGJΦ described in this invention, although other well known methodologies could be applied to remove VGJΦ from an infected strain. Examples of live cholera vaccines to perform the specific mutations mentioned above are vaccine strains which are not able to react with a VGJΦ specific prove and that have been demonstrated in the previous art, to have acceptable levels of reactogenicity in volunteers studies. The genotype of these strains includes suppressive mutations of CTXΦ phage leaving a remnant RS1 and the insertional inactivation of hap gene with the celA gene. Such strains are constructed by means of traditional methods of suppression of the CTXΦ prophage in epidemics strains of V. cholerae , followed by the inactivation of the hemagglutinin protease gene (hap) for the insert of the marker gene celA in their sequence. To see Robert's scientific article and cols., Vaccine, vol 14 No16, 1517-22, 1996), the scientific article of Benítez J. A. and cols, (1999), Infection and Immunity. February; 67(2): 539-45, and the application of invention WO9935271A3, of Campos and cols, 1997. Other strains with these characteristic and that additionally have auxotrophy mutations are also useful to obtain the strains with the characteristics of interest of the present invention. In accordance with the description in the above paragraph, a primary objective of this invention is to protect the use of suppressive or spontaneous mutations conductive to the absence of the MSHA fimbria in the surface of vibrios and in this way impede that live cholera vaccine reacquire and disseminate the cholera toxin genes by means of the infection with the hybrid phage, HybPΦ. Among the preferred inclusions of this invention are any live cholera vaccine strain of the existing biotypes and serotypes or any non toxigenic strain of another emergent serotypes with genetic manipulations that suppress the genome of the CTXΦ phage, inactivate the hap gene, combined with any other mutation, for example the introduction of some auxotrophies (to lysine or metionine) and that also have the characteristics proposed in the present invention. Among the preferred inclusions are also the use of the well tolerated live cholera vaccines, improved by the impossibility of acquiring the CTXΦ in an event mediated for VGJΦ and for the absence of VGJΦ that diminish the risk of dispersion of CTXΦ, as a delivery system to present heterologous antigens to the mucosal immune system. To obtain these mutants in the expression of the MSHA fimbria we have used several molecular biology techniques which are not object of protection of the present document. The present invention also discloses the methods to preserve and lyophilizate these strains with the purpose of being able to prepare live vaccines that present a rapid and adequate reconstitution post lyophilization without affecting their viability when being reconstituted in a solution of sodium bicarbonate 1,33%. It is also the object of the invention that by means of the adequate selection of components, the lyophilized formulations guarantee that live cholera vaccines does not decrease their viability less than 1 logarithmic order as consequence of the storage, independently of the serogroup, serotype or biotype or the mutations they have, even if they were lyophilized for separate or mixed as part of a same preparation. Among the formulation components to be present are lactose (L), peptone (P), yeast extract (E) and sorbitol (S). The total concentration should not exceed the 10%. EXAMPLE 1 Discovery and Characteristics of the VGJφ Phage VGJφ was discovered as an extrachromosomal transmissible element in total DNA preparations from Vibrio cholerae SG25-1, an O139 strain isolated in Calcuta India, 1993 and kindly donated by professor Richard A. Finkelstein. Simple experiments showed the transmissibility of this element. Free-cell culture supernatants of the donor strain, carrying the element, was grown in standard condition like LB media (NaCl 10 g/l, triptone 10 g/l and yeast extract 5 g/l) and was able to transfer to a receptor strain that does not contain any extrachromosomal element, one genetic element of the same size and restriction map that the one was present in the donor strain. The property of transmission without the direct contact between donor-receptor is typical in phages. Infection Assays: Donor strains were grown until optical density to 600 nm equal 0.2. One aliquot from the culture was filtered through a 0.22 μm- pore-size filter to remove the bacteria. The sterility of the filtrate was confirmed by growing one aliquot in LB plates and incubating overnight at 37° C. After checking the lack of colony forming units, 100 μl of free-cell supernatants or serial dilutions were used to infect 20 μl of a fresh culture of the receptor strain. The mixture was incubated for 20 min at room temperature and spread in solid or liquid LB media at 37° C. overnight. The infection was confirmed by the presence of replicative form (FR) and single strand DNA (ssDNA) of VGJφ in the infected vibrios. Purification of VGJφ Phage: Purification of phage particles was done from 100 ml of the culture of 569B Vibrio cholerae strain (classic, Inaba) infected with VGJφ. This strain was used because it contains a CTXφ defective prophage. The cells were centrifuged at 8000×g for 10 min. The supernatant was filtered through a 0.22 μm membrane. Phage particles in the filtrate were precipitated by addition of NaCl and polyethylene glycol 6000 to a final concentration of 3 and 5% respectively. The mixture was incubated in ice for 30 min and centrifuged at 12000×g for 20 min. The supernatant was discarded, and the phage-containing pellet was suspended in 1 ml of phosphate buffered saline. Characterization of VGJφ: VGJφ particles precipitated retained the capacity to infect 569B strain and were stable in PBS solution during at least 6 month at 4° C. After phage particles purification, the genomic DNA was extract using phenol-chloroform solution. The analysis of this DNA showed resistance to digestion with ribonuclease H indicating that the genome is DNA and not RNA (data not shown) and it was also resistant to the treatment with different restriction enzyme but sensitive to treatment with Mung-Bean and Sl nuclease (data not shown), indicating that the phage genome consists of ssDNA. An electrophoresis analysis in the presence of acrydine orange demonstrated similar results to the previous ones. The acrydine orange intercalated in the double stranded DNA (dsDNA) fluoresce green, while fluoresce orange when intercalates in the ssDNA. As expected, the genomic DNA fluoresced orange indicating its single stranded nature (data not shown) and the plasmid DNA observed in the infected cells fluoresced green indicating that it consists of dsDNA. Identity Between the Genome of VGJφ and the Intracellular Replicative Form. Southern blotting analysis carried out using the genome of VGJφ as a probe showed a genetic identity between the extrachromosomal elements of the donor strain SG25-1 and the infected strain 569B. This result confirms that the ssDNA of the viral genome is produced by the cytoplasmic RF and at the same time suggests that VGJφ is a filamentous phage, which uses the rolling circle mechanism of replication to produce the genomic ssDNA that is assembled and exported in phage particles. The RF, isolated from the infected strain 569B, was mapped by restriction analysis. The map obtained showed that the phage genome size (about 7500 b) and the electrophoretic restriction pattern were different to those of the previously reported V. cholerae -specific filamentous phages. These results indicated that the phage isolated from SG25-1 was not described previously and it was designated VGJφ. Titration of VGJφ. For tittering the phage suspensions the procedure was the same as the infection assay, but the indicator strain cells were plated onto an overlay of soft agar (0,4%) over solid LB plates. The plates were incubated overnight at 37° C. and the observed opaque plaques (infection focuses) were counted. This assay revealed that a culture of 569B infected with VGJφ is able to produce until 3×10 11 phage particles per ml of culture, what is unusually high compared with other described filamentous phages of V. cholerae like CTXφ, which produces a maximum of 10 6 particles per ml. Electron Microscopy. Different quantities of VGJφ particles were negatively stained with a solution of 4% uranile acetate (m/v) and observed over a freshly prepared Formvar grids in a transmission electron microscope JEM 200EX (JEOL, Japan). The observation confirmed that the phage particles had a filamentous shape ( FIG. 1 ). Construction and Titration of VGJ-Knφ. The RF of VGJφ was linearized by its unique XbaI site. One DNA fragment containing the R6K replication origin and a kanamycin resistance cassette from pUC4K plasmid was inserted in the XbaI site of VGJφ. This recombinant RF was introduced in V. cholerae 569B and the phage particles were designated as VGJ-Knφ. The donor strain, 569B infected with VGJ-Knφ, was cultured until an OD 600 =2.0. An aliquot of the culture was filtered through a 0.2 um-pore-size filter to eliminate the bacterial cells. The sterility of the cell-free suspension was checked by plating an aliquot of 50 ul in a solid LB plate and incubating overnight at 37° C. Aliquots of 100 ul of the cell-free phage suspension or dilutions of it were used to infect 20 ul of a fresh culture of the receptor strain (about 10 8 cells). The mixture was incubated at RT for 20 min to allow infection. Subsequently, the mixtures were plated onto solid LB supplemented with kanamycin (50 ug/ml) and the plates were incubated overnight at 37° C. The colonies that grow in the presence of antibiotic acquired their Kn-resistance due to the infection with the marked phage VGJ-Knφ. Several of these colonies were checked for the presence of the RF of VGJ-knφ by purification of plasmid DNA and restriction analysis of it. Titration assay done by this method agreed with those obtained by that of opaque plaques with VGJφ, showing that a culture of 569B infected by VGJ-knφ produces about 2×10 11 particles of phage VGJ-knφ per milliter of culture. Nucleotide Sequence: The nucleotide sequence of VGJφ consisted of 7542 nucleotides and had a G+C content of 43.39%. The codified ORFs were identified and compared to protein data bases. The genomic organization of VGJΦ was similar to that of previously characterized filamentous phage, such as phages of Ff group (M13, fd and f1) of E. coli and other filamentous phages of V. cholerae (CTXΦ, fs1, fs2 and VSK) and V. parahemolyticus (Vf12, Vf33 and VfO3k6). VGJΦ does not have a homologous gene to the gene IV of phages of Ff group which suggests that VGJΦ could use a porine of the host for assembling and exporting its phage particles, similar to CTXΦ phage. The nucleotide sequence of VGJΦ revealed that VGJΦ is a close relative of fs1 and VSK phages, sharing several ORF highly homologous and exhibiting 82.8 and 77.8% of DNA homology to VSK and fs1. However, there are genome areas highly divergent and ORFs not share between them. Besides, the genome size is different and it has not been described before that fs1 or VSK being capable of transducing the genes of cholera toxin. The nucleotide sequence of VGJΦ also revealed the presence of two sites homologous to att sequences known to function in integrative filamentous phage. These sites of VGJΦ are partially overlapped and in opposite directions. This arrangement was also found in phages Cf1c, Cf16-v1 and ΦLF of X. campestris as well as Vf33 and VfO3k6 of V. parahemolyticus and VSK of V. cholerae . All these phages except Vf33 and VSK integrate in the chromosome of their hosts by the att site present in the negative strand of the replicative form of these phages. EXAMPLE 2 Identification of VGJΦ Receptor Filamentous phages generally use type IV pili as receptor to infect their hosts. Previously reported V. cholerae -specific filamentous phages use TCP or MSHA pili as receptor. Therefore, two mutants of the El Tor strain C6706 for these pili, KHT52 (ΔtcpA10) and KHT46 (ΔmshA), were used to identify if any of them was the receptor of VGJΦ. While parenteral strain C6706 and its TCP-mutant KHT52 were sensitive to the infection with VGJΦ, the MSHA-mutant KHT46 was fully resistant to the phage, indicating that MSHA was the receptor of VGJΦ. Complementation of strain KHT46 with wild type mshA structural gene (from parental C6706) carried on plasmid pJM132 restored phage sensitivity, confirming that MSHA is the receptor for VGJΦ. The resistance or sensivity to VGJΦ was evaluated by the absence or presence of replicative form in cultures of receptor strain analyzed after the infection assay. To give a numerical titer of the particles which are transduced in each case, it was used an infection assay with VGJΦ-kn as was described previously, resulting the following: The parental strain C6706 and its derivative TCP mutant KHT52 were sensitive to the infection with VGJΦ-Kn and, as indicator strains showed titres of 10 11 plaque forming units (PFU), while KHT46, a MSHA mutant, was fully resistant to the phage, less than 5 PFU/mL, after being infected with the same preparation of VGJΦ-Kn. Complementation of strain KHT46 with wild type msha structural gene, restored phage sensitivity. These results confirm that a mutation that prevents the expression of MSHA pilus confers resistance to the VJGΦ infection. Further assays to compare the capacity of HybPΦ and CTXΦ to infect Clasical and El Tor strains were done, using their kanamycin resistant variants. See the results in Table 1. As it has been previously described, CTXΦ-Kn phage was obtained through the insertion of a kanamycin resistance cassette from the plasmid pUC4K (Amersham Biosciences), in the unique restriction site, NotI, of the replicative form of CTXΦ. The HypPΦ phage was obtained during an infection assay where cell free culture supernatant of 569b strain co-infected with CTXΦ-Kn and VGJΦ-Kn was used to infect the receptor strain KHT52. The cells of this strain carrying kanamycin resistance, originally carried by CTXΦ-Kn and provided to HybPΦ, were purified and, they continued producing HybPΦ viral particles to the supernatant. To check the efficiency of infection of the hybrid phage in Classical and El Tor vibrios, suspensions of CTXΦ-Kn and VGJΦ-Kn of the same title (1-5×10 11 particles/mL) were used to infect the receptor strains 569B (Classical) and C7258 (El Tor). In both cases, the receptor strains were grown in optimal condition for TCP expression, the CTXΦ receptor. The assay was done as follows, 200 μL of pure phage preparation were mix with 20 μL (about 10 8 cells) of a fresh culture of a receptor strain during 20 min at room temperature, plated on solid LB supplemented with kanamycin and incubated over nigh at room temperature. The numbers of colonies carry the Kn-resistance gene in their genome is the result of phage infections and show the capacity of each phage to infect different strains in routine laboratory condition. Those results are exposed in Table 1. TABLE 1 Titration of CTXφ-Kn and hybrid (HybPΦ-Kn) phages in 569B and C7258. Kn r colony number of receptor strain Phage 569B C7258 CTXφ-Kn 5.8 × 10 5 0 HybPΦ-kn 1.5 × 10 5 7.5 × 10 4 As it is shown in Table 1 the hybrid phage transduces CT genes more efficiently than CTXΦ, the ordinary vehicle of these genes. These results point out the importance of the CTXΦ transmission mediated by VGJΦ among Vibrio cholerae strains and stressed its relevance considering the ubiquity of MSHA, the functional receptor in these bacterial strains. EXAMPLE 3 Mobilization of CTXΦ, its Mechanism and Reversion to Virulence Infection of V. cholerae O1 or O139 strains that carry an active CTXΦ phage with VGJΦ gives rise to the production of infective particles that bear the CTXΦ phage genome inserted in the genome of VGJΦ. These particles of hybrid phages have been designated HybPΦ. The HybPΦ titers were evaluated by means of the use of a hybrid phage, which carries a kanamycin marker (HybPΦ-Kn), employing different strains as indicators. The resultant titers are shown in Table 1. HybPΦ-Kn was purified starting from preparations derived of 569B (HybPΦ-Kn) strain and the single strand was sequenced to determine the junctions between CTXΦ and VGJΦ. The cointegrate structure is graphically shown in the FIG. 2 and the nucleotide sequences of the junctions among both sequences, what explains the mechanism by which VGJΦ transduces CTXΦ toward other V. cholerae strains. HybPΦ-Kn enters to V. cholerae using the same receptor that VGJΦ, that is to say MSHA. V. cholerae 1333 strain is an attenuated clone described in the previous art, similar to the strains that were useful for obtaining the derivative of the present invention. This strain is a derivative of the pathogenic C6706 strain. As it shows in the FIG. 4 , the inoculation of 10 5 colony-forming units of 1333 strain in suckling mouse does not have lethal effect, even when it is colonizing for the subsequent 15 days. Several experiments to determine virulence, demonstrated the effect of the HybPΦ-Kn infection on the reversion to virulence. While a dose of 10 5 CFU of 1333 strain does not have a lethal effect, C6706 and 1333 (HybPΦ-Kn) strains have very similar lethality profiles, and don't allow survival of inoculated mouse beyond the fifth day ( FIG. 4 ). EXAMPLE 4 Constitutive Expression of the VGJΦ Receptor, the MSHA Fimbria, in Different Culture Conditions To study the expression of MshA, the major subunit of MSHA fimbria, V. cholerae C7258, C6706, and CA401 strains, were grown in different media. The media used were: LB pH 6.5 (NaCl, 10 g/l; bacteriological triptone, 10 g/l; yeast extract, 5 g/l), AKI (bacteriological peptone, 15 g/l; yeast extract, 4 g/l; NaCl, 0.5 g/l; NaHCO3, 3 g/l), TSB (pancreatic digestion of casein, 17 g/l; papaine digestion of soy seed, 3.0 g/l; NaCl, 5 g/l; dibasic phosphate of potassium, 2.5 g/l; glucose, 2.5 g/l), Dulbecco's (glucose, 4.5 g/l; HEPES, 25 mm; pyridoxine, HCl, HaHCO3), Protein Free Hybridoma Médium (synthetic formulation free of serum and proteins, suplemented with NaHCO3, 2.2 g/l; glutamine, 5 mg/l; red phenol, 20 g/l) and Syncase (NaH2PO4, 5 g/l; KH2PO4, 5 g/l; casaminoacids, 10 g/l; sucrose, 5 g/l and NH4Cl, 1.18 g/l). In all cases was inoculated one colony in 50 ml of culture broth and was grown in a rotary shaker during 16 hours at 37° C., with the exception of the AKI condition in which the strains were grown first at 30° C. in static form during 4 hours and later on rotator shaker at 37° C. during 16 hours. In each case, the bacterial biomass were harvested by centrifugation and used to prepare cellular lisates. Equivalent quantities of cellular lisates were analyzed by Western Blot with the monoclonal antibody 2F12F1 for immunodetection of mshA. The MSHA mutant strain KHT46 was used as negative control of the experiment. All the studied strains, except the KHT46 negative control strain, showed capacity to produce MshA in all culture conditions tested. Equally, said strains cultured in the previous conditions have the capacity to hemagglutinate chicken erythrocytes (mannose sensitive), in the same titer or higher to 1:16 and are efficiently infected by VGJΦ-Kn, exhibiting titers higher than 10 10 particles per milliliter of culture. EXAMPLE 5 Obtaining of Spontaneous Mutants Deficient in MSHA Expression and Evaluation of Resistance to Infection Strain KHT46, a MSHA suppression mutant, derived from V. cholerae C6706 (O1, The Tor, Inaba), shows a refractory state to the infection with VGJΦ, VGJΦ-Kn and the hybrid HybPΦ phages. However, this is a pathogenic strain that is not property of the authors of the present application, neither of the juridical person who presented it, The National Center for Scientific Research, in Havana City, Cuba. To obtain the spontaneous mutants deficient in the expression of superficial MSHA of the present application, was used a suppression mutant in the cholera toxin genes that during the process of obtainment resulted affected in their capacity to assemble MSHA in the cellular surface. Said mutants although are capable of producing the structural subunit of MSHA, do not assemble it in their surface and therefore do not have detectable titers of mannose sensitive hemagglutination, neither adsorb the activity of a specific monoclonal antibody against the MSHA in a competition ELISA. Since this phenotype is notably stable, these mutants were subsequently genetically manipulated to introduce an insertional mutation in the hemagglutinin protease gene, following the procedure described in patent WO 99/35271 “V. cholerae vaccine candidates and the methods of their constructing” of Campos et al, and in the Robert's article, Vaccine, vol 14 No 16, 1517-22, 1996. The resultant mutants were named JCG01 and JCG02, both of O1 serogrup, El Tor biotype, Ogawa serotype. JCG01 and JCG02 showed a refractory state to the infection with the VGJΦ-Kn phage, a variant of the VGJΦ phage that carries a resistance marker to kanamycin. A VGJΦ-Kn suspension that had a proven titer of ˜10 11 units per ml, does not show capacity to infect said strains (non detectable titers, lower to 5 units for ml). This refractory state to the infection with VGJΦ-Kn correspond with a very low titer of hemagglutination in the strains JCG01 and JCG02 (1:2) regarding their parental (1:32) besides a total impairment in the MSHA dependent hemaglutination. Equally, whole cells of these mutants had null capacity to inhibit the interaction of the anti-MSHA monoclonal antibody (2F12F1) to MshA fixed on the solid phase in a competition ELISA. However, both strains produced the major structural subunit MshA, according to immunoblot experiments, indicating that the protein is not correctly assembling in the cellular surface although it is being produced. These mutants allowed proving the concept of this invention and passing to obtain suppression mutants. Obtaining Suppression Mutants in the mshA Gene Starting from Other Cholera Vaccine Candidates. To obtain suppression mutants in the mshA structural gene, two segments of the genome of V. cholerae N16961, of ˜1200 base pairs for each flank of the mshA structural gene were amplified by means of the polimerase chain reaction, using the following oligonucleotides: CNC-8125, ATG ATC GTG AAG TCG ACT ATG (21 mer) (SEQ ID NO:2); CNC-8126 CAG CAA CCG AGA ATT HERE ATC ACC ACG (27 mer) (SEQ ID NO:3); CNC-8127, ATT CTC GGT TGC TGG AAC TGC TTG TG (26 mer) (SEQ ID NO:4); and CNC-8128, GCT CTA GAG TAT TCA CGG TAT TCG (24 mer) (SEQ ID NO:5). The amplified fragments were cloned independently and assembled in vitro to generate the pΔmshA clone. This clone contains these fragments in the same order and orientation that they are found in the bacterial chromosome; only the coding region of the mshA gene has been suppressed from the inner of the sequence. The fragment carrying the suppression was subcloned from the previous plasmid as a Sal I/Xba I fragment in the suicide vector pCVD442 to obtain the plasmid pSΔmshA. The plasmid pSΔmshA was used to suppress the chromosomal mshA gene in the V. cholerae vaccine strains by means of a traditional methodology of allelic replacement. For it, pSΔmshA was introduced in the E. coli strain SM10□pir and mobilized toward V. cholerae by means of a procedure of bacterial conjugation. The resultant clones were selected for their resistance to the ampicillin antibiotic in plates of LB medium supplemented with ampicillin (100 □g/ml) . Most of these clones arise due to integration of the plasmid in the chromosome of the receptor vibrios by means of an event of homologue recombination between one of the flanking fragments to the chromosome mshA gene and that of the plasmid pSΔmshA, originating a cointegrate between both. This event was verified by means of a Southern blot experiment, in which the total DNA of 10 clones was digested with the restriction enzyme Sma I and hybridized with a probe obtained from the plasmid pSΔmshA (Sal I/Xba I insert). The clones of our interest are those that produce a band of 21 000 base pairs. A similar control of the parental strain in this experiment produced a band of 13 000 base pairs. The adequate clones were conserved immediately in LB glycerol at −70° C. Then 3 of them were cultured in the absence of the antibiotic selective pressure to allow that an event of homologue recombination eliminated the genetic duplication existing. This can happen by means of suppression of the original genetic structure (intact mshA gene) and replacement by a mutated copy present in the plasmid (supressed mshA gene) as is shown in FIG. 3 . The clones in which the mutated gene replaced the intact gene were analyzed by Southern blot and identified by the presence of a band of 12 000 base pairs. Finally, the clones where the mshA gene was suppressed were selected and conserved appropriately as vaccine candidates (freezing at −80° C. in LB supplemented with 20% glycerol). This procedure was performed with each clone where the mshA suppression mutant was constructed. Serological Characterization After the introduction of each mutation in the vaccine strains described in this document, each derivative was checked for the correct expression of the lipopolysaccharide corresponding to the original serotype. For that, cells were collected from a fresh plate, resuspended in saline (NaCl, 0.9%) and immediately examined with an appropriate agglutination serum, specific for Ogawa, Inaba or O139 vibrios. The major immune response generated by an anti-cholera vaccine, is against the LPS, therefore the expression of the antigen corresponding to each one of the strains presented in this invention was confirmed by agglutination with specific antiserum. Colonization Assay in Suckling Mice The colonization assay in suckling mice (Herrington et al., J. Exper. Med. 168: 1487-1492, 1988) was used to determine the colonizing ability of each strain. An inoculum of 10 5 -10 6 vibrios in a volume of 50 □l was administered by orogastric route to groups of at least 5 suckling mice. After 18-24 hours at 30° C. the mice were sacrificed, the intestine was extracted and homogenized, and dilutions were plated in appropriate media for the growth of mutants. TABLE 2 Colonizing capacity of the vaccine strains of the present invention. Strain Inoculum Colonizing Genotype BLR01 1.0 × 10 5 2.8 × 10 4 ΔCTXΦ, hap::celA, ΔmshA BLR02 2.0 × 10 6 4.2 × 10 4 ΔVGJΦXΦ, hap::celA, lysA, BLR03 1.2 × 10 6 8.0 × 10 3 ΔCTXΦ, hap::celA, metF, EMG01 3.0 × 10 5 8.0 × 10 6 ΔCTXΦ, hap::celA, ΔmshA EMG02 2.5 × 10 5 3.0 × 10 6 ΔVGJΦXΦ, hap::celA, lysA, EMG03 4.0 × 10 5 5.0 × 10 5 ΔCTXΦ, hap::celA, metF, JCG01 2.0 × 10 5 6.0 × 10 6 ΔCTXΦ, hap::celA, MSHA − JCG02 1.0 × 10 5 6.0 × 10 7 ΔCTXΦ, hap::celA, MSHA − JCG03 1.0 × 10 5 1.0 × 10 6 ΔCTXΦ, hap::celA, ΔmshA EVD01 3.0 × 10 5 3.0 × 10 5 ΔVGJΦXΦ, hap::celA, thyA, KMD01 1.0 × 10 6 7.0 × 10 5 ΔCTXΦ, hap::celA, metF, KMD02 2.0 × 10 6 5.0 × 10 6 ΔCTXΦ, hap::celA, lysA, ESP06 1.7 × 10 6 6.0 × 10 5 ΔCTXΦ, hap::celA, ΔVC0934, JCG04 1.0 × 10 6 2.0 × 10 7 ΔCTXΦ, hap::celA, ΔmshA ESP01 1.0 × 10 5 5.0 × 10 6 ΔVGJΦXΦ, hap::celA, metF, ESP02 6.0 × 10 5 4.0 × 10 5 ΔCTXΦ, hap::celA, lysA, ESP04 8.0 × 10 4 1.0 × 10 6 ΔCTXΦ, hap::celA, ΔVC0934, RAF01 3.1 × 10 5 5.0 × 10 7 ΔCTXΦ, hap::celA, ΔmshA EVD02 2.8 × 10 5 3.1 × 10 6 ΔVGJΦXΦ, hap::celA, thyA, ESP03 1.5 × 10 5 2.0 × 10 6 ΔCTXΦ, hap::celA, metF, KMD03 2.3 × 10 5 3.4 × 10 6 ΔCTXΦ, hap::celA, lysA, ESP05 2.1 × 10 6 2.3 × 10 6 ΔCTXΦ, hap::celA, ΔVC0934, TLP01 2.3 × 10 6 3.2 × 10 5 ΔCTXΦ, hap::celA, ΔmshA TLP02 3.4 × 10 5 9.4 × 10 4 ΔVGJΦXΦ, hap::celA, lysA, TLP03 2.7 × 10 5 8.8 × 10 4 ΔCTXΦ, hap::celA, metF, All the strains showed adequate colonizing capacity to be used as live vaccine candidates. The colonization is needed to generate a strong immunological response because the local multiplication of the bacteria increases the duration of interaction with the mucosal immune system. In this case, although a perfect model for cholera does not exist, the suckling mice gives an adequate approach to what can be the subsequent colonization of each strain in humans. Motility Assay The cells of a well isolated colony are loaded in the tip of a platinum loop from a master plate toward a plate for the motility detection (LB, agar 0.4%), introducing the tip of the loop 2-3 mm in the agar. The diameter of dispersion of each colony in the soft agar to 30° C. is measured at 24 hours of incubation. A bacterial strain that reaches a diameter of 3 mm or less from the point of inoculation is considered as non-motile. A bacterial strain that grows in a diameter beyond 3 mm is considered as motile. All the strains included in this invention resulted to be motile. EXAMPLE 6 Methods to Select and Construct the Vaccine Candidates Useful as Starting Strains to be Modified by the Procedure Disclosed in the Present Invention Five pathogenic strains in our collection were selected as starting microorganism due to their lack of hybridization with VGJΦ sequences. These strains are V. cholerae C7258 (O1, El Tor, Ogawa, Perú, 1991), C6706 (O1, El Tor, Inaba, Perú, 1991), CRC266 (O139, La India, 1999), CA385 (Clásico, Ogawa) y CA401 (Clásico, Inaba). The procedures disclosed in this example are not the subject of the present invention. They rather constitute a detailed description of the methods used to obtain attenuated strains that are the substrate to construct the mutants claimed in the present invention. These mutants being characterized in that they are refractory to infection by VGJΦ and the hybrid VGJΦ::CTXΦ are obtained by the methods described in the examples 4 and 5. Below we describe the suicide plasmids used to introduce different sets of mutations into V. cholerae by allelic replacement before they are suitable to be modified by the methods of the present invention. The reader should note that the strains claimed in the present invention have in addition to the mutation that impairs the correct expression of MSHA fimbriae (a) a deletion mutation of the cholera enterotoxin genes or the entire CTXΦ prophage and (b) the hemaglutinin protease gene interrupted with the Clostridium thermocellum endoglucanase A gene. They can also have additionally and optionally mutations in the genes (c) lysA, (d) metF, (e) VC0934 (coding for a glycosil transferase) and (f) thyA. (a) To construct atoxigenic strains by inactivation of the cholera enterotoxin genes or deletion of the CTXΦ prophage, the suicide plasmid used was pJAF (Benitez y cols, 1996, Archives of Medical Research, Vol 27, No 3, pp. 275-283). This plasmid was obtained from plasmid pBB6 (Baudry y cols, 1991, Infection and Immunity 60:428), which contains a 5,1 kb insert from V. cholerae 569B that encodes ace, zot, ctxA y ctxB. Due to the absence of RS1 sequences 3′ to the ctxAB operon in Classical vibrios, the EcoR I site downstream to the ctxAB copy in this plasmid lies in the flanking DNA of undefined function. The plasmid pBB6 was modified by deletion of the ScaI internal fragment to create plasmid pBSCT5, which now contains a recombinant region deprived of the zot and ctxA functional genes. Then the PstI of pBSCT5 was mutated into EcoRI by insertion of an EcoRI linker to obtain pBSCT64 and the resultant EcoRI fragment was subcloned into the EcoRI site of pGP704 to obtain pAJF. (b) To construct strains affected in the expression of HA/P the suicide plasmid pGPH6 was used. This plasmid was constructed in different steps. First, plasmid pCH2 (Hase y Finkelstein, 1991, J. Bacteriology 173:3311-3317) that contains the hap gene in a 3,2 kb HindIII fragment from V. cholerae 3083 was linealized by the StuI site, which is situated in the hap coding sequence. The 3.2 kb HindIII-fragment containing the celA gene was excised from plasmid pCT104 (Cornet y cols, 1983, Biotechnology 1:589-594) and subcloned into the StuI site of pCH2 to obtain pAHC3. The insert containing of pAHC3, containing the hap gene insertionally inactivated with the celA gene, was subcloned as a HindIII fragment to a pUC19 derivative that have the multiple cloning site flanked by BglII sites to obtain pIJHCI. The Bgl II fragment of this plasmid was subcloned into pGP704 to originate pGPH6, which contains a 6.4 kb fragment with the genetic hap::celA structure, where the hap gene is not functional. (c) y (d) When constructing mutants in the lysA or metF genes, the suicide plasmids pCVlysAΔ1 or pCVMΔClaI were used. To construct these plasmids, the lysA y metF genes were PCR amplified from V. cholerae C7258, using a pair of oligonucleotides for each gene. The oligonucleotides were purchased from Centro de Ingeniería Génetica y Biotecnología, Ciudad de La Habana, Cuba. The nucleotide sequences of the primers were: (lysA): (P 6488) 5′-GTA AAT CAC GCT ACT AAG-3′(SEQ ID NO:11) and (P 6487) 5′-AGA AAA ATG GAA ATGC-3′(SEQ ID NO:10) and (metF): (P 5872) 5′-AGA GCA TGC GGC ATG GC-3′(SEQ ID NO:8) and (P 5873) 5′-ATA CTG CAG CTC GTC GAA ATG GCG-3′(SEQ ID NO:9). The amplicons were cloned into the plasmids pGEM®T (Promega) and pIJ2925 (Janssen y cols, 1993, Gene 124:133-134), leading to the obtainment of the recombinant plasmids pGlysA3 y pMF29, which contain active copies of the lysA y metF genes, respectively. The identity of each gene was checked by nucleotide sequencing. The metF and lysA genes cloned were mutated in vitro by deletion of the respective ClaI (246 base pairs) and PstI/AccI (106 base pairs) inner fragments, respectively. In the last of the cases the strategy was designed to keep the open reading frame leading to an inactive gene product to avoid exerting polar effects during and after construction of a lysA mutant of V. cholerae . Each inactivated gen was cloned as a Bgl II fragment in the suicide vector pCVD442 for the subsequent introduction into the cholera vaccine candidates of interest. The suicide plasmid containing the lysA alelle was termed pCVlysAΔ1 and the one containing the metF alelle was denominated pCVMΔClaI. (e) When constructing mutants of the VC0934 gene we constructed and used the suicide plasmid pCVDΔ34. In doing that, the VC0934 gene was PCR amplified using as template total DNA from strain N16961 and the primers: 5′-GCA TGC GTC TAG TGA TGA AGG-3′(SEQ ID NO:6) and 5′-TCT AGA CTG TCT TAA TAC GC-3′(SEQ ID NO:7) The amplicon was cloned into the plasmid pGEM5Zf T-vector to obtain plasmid pGEM34; a 270 base pair deletion was performed inside the VC0934 coding sequence using the restriction enzymes NarI/BglII. After flushing the ends with klenow and subsequent recircularization the plasmid obtained was named pG34. The resultant inactive gene was subcloned into the suicide vector pCVD442 digested with SalI and SphI to obtain the plasmid pCVΔ34. This plasmid was used to make the allelic replacement of the wild type gene. (f) When constructing mutants defective in thyA expression the suicide plasmid pEST was constructed and used. The steps to construct this plasmid comprised the cloning of the thyA gene from V. cholerae C7258 into pBR322 as an EcoRI-HindIII cromosomal DNA fragment, to obtain pVT1 (Valle y cols, 2000, Infection and Immunity 68, No 11, pp6411-6418). A 300 base pairs internal fragment from the thyA gene, comprised between the BglII and MluI, sites was deleted from this plasmid to obtain pVMT1. This deletion removed the DNA fragment that codes for amino acids 7 to 105 of the encoded protein Thimidilate syntase. The mutated thyA gene was excised as an EcoRI-HindIII fragment, the extremes were blunted and then cloned into the SmaI site of pUC19 in the same orientation as the β-galactosidase gene to obtain pVT9. The resultant gene was subcloned as a SacI fragment from pVT9 to pCVD442 and the obtained plasmid was named pEST. This final construct was used to make the allelic replacement of the wild type gene in the strains of interest. The described suicide vectors are a modular system that can be used to introduce secuencial mutation into V. cholerae vaccine candidates. The allelic replacement with the genes encoded by these vectors is done following the sequence of steps denoted below: In the first step the suicide vector, containing the allele of choice among those described, is transferred by conjugation from the E. coli donor SM10λpir to the V. cholerae recipient, this last being the subject of the planned modification. This event is done to produce a cointegrate resistant to ampicillin. The clones resultant from the conjugational event are thus selected in LB plates supplemented with ampicillin (100 μg/ml). The procedure for this first stage is as follows. The donor strain, SM10λpir transformed with the sucide vector of interest, is grown in an LB plate (NaCl, 10 g/l; bacteriological triptone, 10 g/l, and yeast extract, 5 g/l), supplemented with ampicillin (100 zg/ml), and the receptor strain, the V. cholerae strain to be modified, is grown in an LB plate. The conditions for growth are 37° C. overnight. A single colony of the donor and one from the receptor is streaked into a new LB plate. The donor strain is streaked firstly in one direction and the receptor ( V. cholerae ) secondly in the opposed orientation. This perpendicular and superimposed streaking warrant that both strain grow in close contact. In the next step the plates are incubated at 37° C. for 12 hours, harvested in 5 ml of NaCl (0.9 %) y 200 μl of dilutions 10 2 , 10 3 , 10 4 and 10 5 are disseminated in LB-ampicillin-polimixinB plates, to select the V. cholerae clones that were transformed with the suicide plasmid and counterselect the donor E. coli SM10λpir. Ten such clones resulting from each process are preserved frozen at −80° C. in LB-glicerol at 20% to be analyzed in the second step. In the second step, a Southern blot hybridization is performed with a probe specific for the gene subjected to the mutational process; this is done to detect the structure of the correct cointegrate among the clones conserved in the previous step. The clones in which the suicide plasmid integrated to the correct target by homologous recombination are identified by the presence of a particular cromosomal structure. This structure contains one copy of the wild type gene and one of the mutated allele separated only by plasmid vector sequences. This particular structure produces a specific hybridization pattern in Southern blot with the specific probe that allows its identification. The appropriate clones are conserved frozen at −80° C. in LB glicerol. This second step comprises the following substeps: Firstly, the total DNA of each clone obtained in the first step is isolated according to a traditional procedure (Ausubel y cols, Short protocols in Molecular Biology, third edition, 1992, unit 2.4, page 2-11, basic protocol). Total DNA from the progenitor strain is isolated as control. Then, the total DNA of the ten clones the progenitor strain is digested with the appropriate restriction enzymes, to be mentioned subsequently in the document. One μg of DNA are digested from each clone and the mother strain and later electrophoresed in parallel lanes of an agarose gel. The DNA content of the gel is blotted into membranes in alkaline conditions (Ausubel y cols, Short protocols in Molecular Biology, third edition, 1992, unit 2.9 A, page 2-30, alternate protocol 1). The blots are fixed by incubation at 80° C. for 15 minutes. The free sites in the membrane are then blocked by prehybridization and subsequently probed with the specific probe for each mutation. What follows are the details of the restriction enzyme, the probe (digoxigenin-labelled using the method random primed method) and the size of the hibridization fragment that identify the desired structure for the cointegrate of each clone, according to the target gene: For suppression mutants of the CTXΦ phage genes, the total DNA of clones is digested with the restriction enzyme Hind III, and once in the membrane is hybridized with a probe obtained starting from the Pst I-EcoR I fragment of the pBB6 plasmid. The clones of interest are the ones that have the genetic structure that origin two bands in the Southern blot, one of 10 000 base pairs and another of 7 000 base pairs. As control the parental strain origins a single band of 17 000 base pairs in the same experiment of Southern blot. For suppression mutants of the hap gene, the total DNA of clones is digested with the restriction enzyme Xho I, and once in the membrane is hybridized with a probe obtained starting from the Hind III fragment of 3 200 base pairs presents in the pCH2 plasmid. The clones of interest are those that have the genetic structure that origins a single band in the Southern blot, of 16 000 base pairs. As control the parental strain generates a single band of 6 000 base pairs in the same experiment of Southern blot. For suppression mutants of lysA gene, the total DNA of clones is digested with the restriction enzyme Xho I, and once in the membrane is hybridized with a probe obtained from the Sph I/Sma I fragment of the pCV□lysAl plasmid, contained the mutated gene lysA. The clones of interest are those that have the genetic structure that origins a single band in the Southern blot, of 12 500 base pairs. As control the parental strain generates a single band of 5 200 base pairs in the same experiment of Southern blot. For suppression mutants of metF gene, the total DNA of clones is digested with the restriction enzyme Nco I, and once in the membrane is hybridized with a probe obtained from the Bgl II fragment of pCVM□ClaI, contained the mutated metF gene. The clones of interest are those that have the genetic structure that origins a single band in the Southern blot, of 12 000 base pairs. As control the parental strain generates a single band of 5 000 base pairs in the same experiment of Southern blot. For suppression mutants of gene VC0934, the total DNA of clones is digested with the restriction enzyme Ava I, and once in the membrane is hybridized with a probe obtained from the Sal I/Sph I fragment of pCVD□34, contained the mutated VC0934 gene. The clones of interest are those that have the genetic structure that origins two bands in the Southern blot, one of 1 600 or 1 900 and another of 8 200 or 7 900 base pairs. As control the parental strain generates a single band of 3 500 base pairs in the same experiment of Southern. For suppression mutants in thyA gene, the total DNA of clones is digested with the restriction enzyme Bstx I, and once in the membrane is hybridized with a probe obtained from the Sac I fragment of pEST1, contained the mutated thyA gene. The clones of interest are those that have the genetic structure that origins a single band in the Southern blot, of 9 600 base pairs. As control the parental strain generates a single band of 2 400 base pairs in the same experiment of Southern blot. In the third step of the procedure, 3 clones of interest, carrying a cointegrate with one of the previous structures, are cultured in absence of the antibiotic selective pressure to allow the loss of the suicidal vector by means of homologue recombination and the amplification of resultants clones. In said clones the loss of the suicidal vector goes with the loss of one of the two copies of the gene, the mutated or the wild one, of the genetic endowment of the bacteria. In a fourth step of the procedure, dilutions of the previous cultures are extended in plates to obtain isolated colonies, which are then replicated toward plates supplemented with ampicillin to evaluate which clones are sensitive to ampicillin. Said clones, sensitive to ampcillin, are conserved for freezing, as described previously. In a fifth step, by means of a study of Southern blot with specific probes for each one of the genes of interest (describe in a, b, c, d, and, f) it is verified which clones retained in the chromosome the mutated copy of the allele of interest. These clones of interest are expanded to create a work bank and to carry out their later characterization, as well as the introduction of the modifications object of protection in the present invention application. In the following paragraphs we detail the restriction enzyme, the probe and the sizes of the hybridization fragments that identify the desired structure in each of the mutants, according to each of the genes being the subject of modification: To analyze the mutants in the CTXΦ prophage, the total DNA is digested with the restriction endonuclease Hind III. Once in the membrane it is hybridized with a probe derived from the Pst I-EcoR I fragment of plasmid pBB6. Are clones of interest such that do not produce hybridization bands in the Southern blot. For the mutants with the inactivated allele of hap, total DNA from the clones is digested with the restriction enzyme Xho I and once in the membrane it is hybridized with a probe derived from the 3 200 base pair Hind III fragment from plasmid pCH 2 that codes for the hap gen. The clones of interest are those that produce a single band in the Southern blot, of about 9 000 nucleotide pairs. For the mutants in the lysA gene total ADN is digested with the restriction enzyme Xho I, and once in the membrane it is hybridized with a probe derived from the Sph I/Sma I fragment isolated from the plasmid pCVΔlysA, that contain the lysA mutated gene. The clones of interest are those having the genetic structure that produce a single band in Southern blot of about de 5 000 pairs of nucleotides. For the mutants with deletions in the metF gene, the total ADN of the clones is digested with the restriction enzyme Nco I, and once in the membrane it is hybridized with a probe derived from the Bgl II fragment contained in plasmid pCVMΔClaI, that contains the metF mutant gene. The clones of interest are those that have the genetic structure that leads to a single band of 4 700 base pairs in the Southern blot. For the mutants in the VC0934 gene, total DNA of the clones is digested with the restriction enzyme Ava I, and the blots are hybridized with a probe obtained from the Sal I/Sph I fragment of pCVDΔ34, which contains the VC0934 mutant gene obtained in vitro. The clones of interest are those having the structure leading to a single band of 3 200 base pairs in the Southern blot. For the mutants in the thyA gene, total DNA of the clones is digested with the restriction enzyme Bstx I, and the blots are hybridized with a probe obtained from the Sac I fragment of pEST1, which contain the thyA gene. The clones of interest are those that have the genetic structure leading to a single band in the Southern blot of about 2 100 base pairs. EXAMPLE 7 Methods to Preserve Vaccine Strains by Means of Lyophilization In following example microorganisms were cultured in LB broth at 37° C. with an orbital shaking 150 and 250 rpm until reaching the logarithmic phase. Cells were harvested by centrifugation 5000 and 8000 rpm at 4° C. during 10-20 minutes and then were mixed with the formulations that show good protection features of the microorganism, so that the cellular concentration was between 10 8 and 10 9 cells ml −1 . 2 ml were dispensed for each 10R type flask. The lyophilization cycle comprised a deep freezing of the material, a primary drying keeping each product between −30° C. and −39° C. for space of 8 to 12 hours and a secondary drying at temperatures between 18° C. and 25° C. for not more than 12 hours. The viability loss was defined as the logarithmic difference of the CFU/mL before and after the lyophilization or before and after the storage of the lyophilized material, which is always dissolved in a 1.33% sodium bicarbonate solution. Formulation L+E+S The BLR01, JCG03 and ESP05 strains were processed by the previously described lyophilization process in a formulation of the type L (5.0%), E (2.0%) and S (2.0%). The freezing was performed at −60° C. During the primary drying, the temperature of the product was kept at −32° C. for 10 hours and in the secondary drying the temperature was kept at 22° C. for 12 hours. The dissolution of the lyophilized material in a 1.33% sodium bicarbonate solution was instant. The viability loss calculated immediately after the dissolution, with regard to the concentration of live cells before the lyophilization resulted to be 0.30, 0.43, and 0.60 logarithmic orders for BLR01, JCG03 and ESP05, respectively. Comparison of the L+P+S and L+E+S Formulations with that of Skim Milk+Peptone +Sorbitol The strain JCG03 was lyophilized using two formulations: the type L (6.0%), P (2.0%) and S (2.0%), and the other type L (5.5%), E (1.8%) and S (1.6%). This strain was also lyophilized in a formulation of 6.0% skim milk, 2.0% peptone and 2.0% sorbitol as a comparison formulation. The freezing was done at −60° C. During the primary drying, the temperature of the product was kept at −33° C. for 12 hours and in the secondary drying the temperature was kept at 20° C. for 14 hours. The dissolution of the lyophilized material in a 1.33% sodium bicarbonate solution was instant when the lyophilization process took place in the formulations of the type L+P+S or L+E+S and slightly slower when was lyophilized in the comparison formulation. The viability loss calculated immediately after the dissolution, with regard to the concentration of live cells before the lyophilization resulted to be 0.48, 0.52 and 0.55 logarithmic orders for the L+P+S, L+E+S and the comparison formulations, respectively, significantly similar. Humidity and Oxygen Effects The strain JCG03 lyophilized in the three formulations mentioned in the previous paragraph, was exposed immediately after being lyophilized to the simultaneous action of humidity and oxygen. This was achieved, confining the samples during 3 days at 25° C. in an atmosphere in sterile glass desiccators, under an 11% relative humidity (created by a saturated solution of lithium chloride). The viability loss in the L+P+S, L+E+S and comparison formulations resulted to be 1.61, 1.10 and 3.43 logarithmic orders, respectively, what shows that the formulations object of this invention guarantee a bigger protection to humidity and oxygen than the comparison formulation. Effect of the Storage Temperature The strains TLP01, JCG01 and ESP05 were lyophilized in a formulation of the type L (5.5%), E(2.0%) and S(2.0%). The freezing was done at −58° C. During the primary drying, the temperature of the product was kept at −30° C. for 12 hours and in the secondary drying the temperature was kept at 20° C. for 14 hours. The dissolution in a 1.33% sodium bicarbonate solution was instant. The viability loss calculated immediately after the dissolution, with regard to the concentration of live cells before the lyophilization resulted to be 0.43, 0.55 and 0.44 logarithmic orders in TLP01, JCG01 and ESP05, respectively. The lyophilized material was stored 1 year either at 8° C. or −20° C. The Table 3 shows the viability loss results obtained. TABLE 3 Viability loss (1 year of storage). Strain 8° C. −20° C. TLP01 1.07 0.64 JCG01 1.02 0.59 ESP05 0.91 0.55 EXAMPLE 8 Strains of the Present Invention and Their Characteristics The strains of the present invention have been deposited on Dec. 11, 2003 in the Belgium Coordinated Collection of Microorganisms (BCCM), Laboratorium voor Microbiologie-Bacterienverzameling (LMG), Universiteit Gent, K. L. Ledeganckstraat 35, B-9000 Gent, Beigium: Vibrio cholerae JCG01 (LMG P-22149) Vibrio cholerae JCG02 (LMG P-22150) Vibrio cholerae JCG03 (LMG P-22151) Vibrio cholerae KMD01 (LMG P-22153) Vibrio cholerae KMD02 (LMG P-22154) Vibrio cholerae KMD03 (LMG P-22155) Vibrio cholerae JCG04 (LMG P-22152) Vibrio cholerae ESP01 (LMG P-22156) Vibrio cholerae ESP02 (LMG P-22157) Vibrio cholerae ESP03 (LMG P-22158) Vibrio cholerae RAF01 (LMG P-22159) Vibrio cholerae TLP01 (LMG P-22160) Vibrio cholerae TLP02 (LMG P-22161) y Vibrio cholerae TLP03 (LMG P-22162) They are described in Table 4. TABLE 4 Vaccine strains of the present invention. Wild type parental Biotype/ Strain strain Serotype Relevante Genotype. BLR01 CA385 Classical/Og ΔCTXΦ, hap::celA, ΔmshA BLR02 CA385 Classical/Og ΔCTXΦ, hap::celA, lysA, ΔmshA BLR03 CA385 Classical/Og ΔCTXΦ, hap::celA, metF, ΔmshA EMG01 CA401 Classical/In ΔCTXΦ, hap::celA, ΔmshA EMG02 CA401 Classical/In ΔCTXΦ, hap::celA, lysA, ΔmshA EMG03 CA401 Classical/In ΔCTXΦ, hap::celA, metF, ΔmshA JCG01 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, MSHA − JCG02 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, MSHA − JCG03 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, ΔmshA EVD01 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, thyA, ΔmshA KMD01 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, metF, ΔmshA KMD02 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, lysA, ΔmshA ESP06 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, ΔVC0934, JCG04 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, ΔmshA ESP01 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, metF, ΔmshA ESP02 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, lysA, ΔmshA ESP04 C7258 El Tor/Ogawa ΔCTXΦ, hap::celA, ΔVC0934, RAF01 C6706 El Tor/Inaba ΔCTXΦ, hap::celA, ΔmshA EVD02 C6706 El Tor/Inaba ΔCTXΦ, hap::celA, thyA, ΔmshA ESP03 C6706 El Tor/Inaba ΔCTXΦ, hap::celA, metF, ΔmshA KMD03 C6706 El Tor/Inaba ΔCTXΦ, hap::celA, lysA, ΔmshA ESP05 C6706 El Tor/Inaba ΔCTXΦ, hap::celA, ΔVC0934, TLP01 CRC266 O139 ΔCTXΦ, hap::celA, ΔmshA TLP02 CRC266 O139 ΔCTXΦ, hap::celA, lysA, ΔmshA TLP03 CRC266 O139 ΔCTXΦ, hap::celA, metF, ΔmshA Advantages The present invention provide us with a methodology to protect live cholera vaccine candidates from the reacquisition of cholera toxin genes and others toxins from the CTXΦ bacteriophage mediated by VGJΦ phage, and therefore from the conversion to virulence by this mechanism. Equally provide us with the necessary information to assure that live cholera vaccine candidates will not spread CTXΦ, in the case that these vaccine candidates reacquire CTXΦ, by a specialized transduction with the VGJΦ phage. The present invention provide us the application of MSHA mutants as live cholera vaccine candidates, which exhibits an increase in their environmental safety due to resistance to the infection with CTXΦ mediated by VGJΦ. This invention provides us with a new characteristic to keep in mind during the design and construction of live cholera vaccine candidates to improve their environmental safety, that is to say that such vaccines are not able to spread the CTXΦ genes, mediated by VGJΦ, in the case of reacquisition. The above characteristic could be applied to the already made live cholera vaccine candidates, which have demonstrated an acceptable level of reactogenicity in volunteers studies, to reduce their potential environmental impact. This invention also provide formulations to preserve by lyophilization all of the above-mentioned live cholera vaccine candidates and also improve their abilities to tolerate the remainder of oxygen and humidity in the container. These formulations also guarantee the instant reconstitution of the lyophilized live cholera vaccine candidate powder in sodium bicarbonate buffer, making easier the manipulation, protecting the vaccines during this process and improving the organoleptic characteristic, specifically related with the visual aspect of lyophilized tablets and the reconstituted products. One of the formulations provided here for conservation and lyophilization of live cholera vaccine candidates lack the bovine components usually added to many formulations to lyophilize human vaccines.

The present invention discloses new live attenuated strains for oral immunization against cholera that are provided in freeze dried formulations for long term storage and administration to humans. These strains combine the two most important properties of live attenuated cholera vaccine candidates. One such property is being well tolerated by people ingesting them. This was achieved by virtue of mutations already described in the art. The second property is having enhanced environmental safety due to the absence of VGJΦ DNA in their genomes and also due to null mutations in the mshA gene or other spontaneous mutations conducive to the lack of MSHA type IV fimbria on the bacterial surface. This was done envisioning that VGJΦ is a filamentous phage able to recombine with CTXΦ and disseminate the cholera toxin genes. This VGJΦ phage as well as the VGJΦ-CTXΦ recombinants uses the MSHA fibers as receptor. Being devoid of MSHA fimbria the vaccine candidates are protected from acquiring CTXΦ from the recombinant hybrid VGJΦ-CTXΦ. Being devoid of VGJΦ, the vaccine candidates are impaired in the dissemination of CTXΦ, via VGJΦ.

BACKGROUND OF THE INVENTION 1. Field of the Invention This invention relates to an apparatus and method for measuring the resistance of superconductors. Structures of this type, generally, allow the resistance of the superconductor to be accurately measured in a non-destructive manner by using a bifilar coil which includes an integrated loop/switch formed from the bifilar coil. 2. Description of the Related Art The use of superconductors for production devices has a problem in that the resistance (or absence thereof) of a long length of material cannot be easily measured nondestructively, i.e. without actually building the device and measuring its field accurately over time. The reason for this is that in producing the desired field, the current in the conductor also produces large forces on the conductor which require the conductor to be potted in a hard material, such as epoxy to avoid motion and consequent quenching of the superconductor. One way to address this issue is to wind a bifilar coil, wherein two conductors are wound adjacent to one another so that current will flow in opposite directions when the winding is excited. The field thus produced is very low and the coil can typically be wound "dry" without concern about the forces, the resultant motion and quenching. This facilitates unwinding after a successful test and subsequent rewinding into the desired field-producing coil configuration. The bifilar coil can then be excited from its end leads and the voltage across the leads measured at a given current to infer the resistance of the conductor. The measurement which is desired from such a screening test is one of the overall resistance of the length of superconductor together with any joints which may be present in it. Unfortunately, the resistance which is required for a persistent current device such as an MR magnet is very low, and therefore difficult to detect by voltage measurements. An estimate of the resistance level which must be attained for a given device can be made from an equation governing the behavior of a circuit with a series resistance (R) and inductance (L) (FIG. 1). For a given initial current I o at time 0, the current at time t is given by I=I.sub.o e.sup.-tR/L (Eq. 1). For example, the specification for field drift rate for an MR magnet is typically 0.1 ppm/hr. Solving Equation 1 for R/L with this desired current drift rate (the current and field are linearly related) yields R/L=2.8e -11 Ω/H. For a magnet with an inductance of 20 Henries, the overall coil resistance must be no greater than 5.6e -10 Ω. At a current of 150 A, the voltage generated by such a resistance is 83e -9 V. Because of thermally induced voltages in the sensing leads and at junctions, as well as, other effects, the accurate measurement of such a low voltage is very difficult-a typical limitation on accurate voltage measurement is perhaps 100e -9 V. Therefore, inference of resistance of the sample of superconductor from the measured voltage can lead to results which are not of the required accuracy. Resistance tests of superconductors have been made in the past by measuring the voltage at a given current for long (up to 70,000 feet) and short (4-12 inches) lengths of conductor. Resistance has been measured from drift tests on individual loops of conductor, but this qualifies only that short length of material and is not useful for qualifying lengths which are required to wind a complete magnet. Therefore, a more advantageous system, then would be presented if the resistance of long lengths of superconductor could be measured in a non-destructive manner. It is apparent from the above that there exists a need in the art for a system which measures the resistance of a superconductor, and which at least equals the voltage measurement accuracy of the known measurement systems, but which at the same time is capable of measuring long lengths of superconductor. It is a purpose of this invention to fulfill this and other needs in the art in a manner more apparent to the skilled artisan once given the following disclosure. SUMMARY OF THE INVENTION Generally speaking, this invention fulfills these needs by providing an apparatus for measuring the resistance of a superconductor, comprising a bifilar coil having a loop formed substantially integrally with said coil and a field measuring means located substantially adjacent to said loop. In certain preferred embodiments, the bifilar coil is wound with paired conductors which carry current in opposite directions to minimize the field, and the forces created in the coil. The coil includes a loop which is formed from the radially outboard of the two paired conductors in the last (outboard) layer. The loop performs two functions with a heater attached. The loop acts as a switch to allow the bifilar coil to be ramped up to a desired current level, and the loop also acts as a field generator which can be used to infer the current in the conductor. The drift rate of the current, coupled with the known (calculated or measured) inductance of the bifilar coil, allows an accurate computation of the overall resistance of the coil and any joints therein. A direct extension of the proposed concept performs the resistance measurement in a background magnetic field. The advantage of this method is that it allows exploration of the resistance of the tested conductor in the actual operating regime which it will see the when the conductor is wound into a magnet. This magnetic field may be imposed on the coil by a solenoidal or other suitable winding which is powered by a highly stable power supply, or which is separately rendered persistent with a superconducting switch. Background magnetic field effects on the desired field measurement are minimized by ensuring that the background field is orthogonal to the field of the loop and by using field measurement techniques which integrate the field around a loop which contains no part of the field coil. In another further preferred embodiment, the resistance of the superconductor can be accurately measured in a non-destructive manner. The preferred superconductor measurement system, according to this invention, offers the following advantages: ease of use; excellent resistance measurements characteristics; excellent economy; good stability; and high strength for safety. In fact, in many of the preferred embodiments, these factors of use, resistance measurement characteristics, and economy are optimized to an extent that is considerably higher than heretofore achieved in prior, known superconductor measurement systems. BRIEF DESCRIPTION OF THE INVENTION The above and other features of the present invention which will become more apparent as the description proceeds are best understood by considering the following detailed description in conjunction with the accompanying drawings wherein like characters represent like parts throughout the several views and in which: FIG. 1 is a R-L circuit diagram for a superconducting coil, according to the prior art; FIG. 2 is a circuit diagram for the switch/coil, according to the present invention; FIG. 3 is a schematic illustration of the loop/switch for a bifilar winding, according to the present invention; FIG. 4 is a second embodiment of the loop/switch for a bifilar winding; FIG. 5 is a third embodiment of the loop/switch for a bifilar winding; FIG. 6a is a schematic illustration of a apparatus for placing a background magnet field on a bifilar coil, according to the present invention; FIG. 6b is a second embodiment of an apparatus for placing a background magnetic field on a bifilar coil; and FIG. 6c is a third embodiment of an apparatus for placing a background magnetic field on a bifilar and a loop/switch for the bifilar coil. DETAILED DESCRIPTION OF THE INVENTION As discussed earlier, FIG. 1 is a prior art diagram for a R-L circuit 2 for a conventional superconducting coil. Circuit 2 includes power source (V) 4, resistor (R) 6 and inductor (L) 8. As disclosed by the present invention, a more accurate measurement of the resistance of a length of conductor than is obtainable from voltage measurements can be obtained by connecting the conductor with superconducting joints in a closed loop circuit, inducing a current into it, and reading the field created by its current over a long time. Since the field varies linearly with the current, this amounts to measuring the current drift. The resistance of a length of conductor with an inductance L over a time interval Δt is given by R=(0.05)L/Δt (Eq. 2). A reasonable time interval over which to measure the field is one day, or 86,400 s. A bifilar coil typically wound of 70,000 feet of tape superconductor (enough to make a 0.5T MR magnet) has an inductance of about 1200 μH. This value is independent of the diameter and length of the coil, as it involves only the spacing between and the length of the two conductors. It may be derived as the inductance of a pair of parallel conductors carrying oppositely directed currents from Gauss' Law in the form: ##EQU1## where λ is the inductance in Henries/meter, g and w are the gap between and width of the conductors, respectively, and μ o is the magnetic permeability of free space. Alternatively, the inductance may be experimentally derived from measurement of the voltage required to create a given rate of current change in the coil. Such a measurement should be carried out over a region of no appreciable flux penetration into the superconductor, or preferably in the normal state but at low temperature so the resistive effects are minimized. Using the 1200 μH calculated value of inductance, from Eq. 3, the resistance from Eq. 2 corresponding to a 4% field drop over one day is 5.6e-10Ω. This is equal to the resistance value which must be measured. The 4% field measurement accuracy is well within the range of a Hall sensor. This preferred test can be accomplished with a superconductor by attaching the current leads to form a parallel circuit having as one leg a short length of conductor (called a switch) which may be driven above its transition temperature with a small heater and as the other the bulk of the coil (see FIG. 2). In particular, FIG. 2 illustrates circuit 10 including power source 12, switch resistor 14 and coil resistor 16. With respect to FIG. 3, the simplest embodiment of switch assembly 20 is formed by wrapping along the direction of arrow A a portion of the radially outermost turn of one conductor 24 of the paired bifilar conductors 22 around a mandrel 26 of the desired diameter in the opposite direction to that of the plural turn coil winding 21 formed by said bifilar conductors which is being wrapped in the direction of arrow B for a single turn to form switch 28 on mandrel 26 before it rejoins the radially outermost turn of the other conductor 30 to complete the final turn of the bifilar coil with a pigtail-type joint 32. In order to induce the desired current in coil winding 21, switch 28 is turned normal by powering a conventional heater 34 attached to switch 28 by conventional fasteners, energizing power leads 36, and allowing the current to stabilize. The currents flowing in the coil (I c ) and the switch (I s ) are determined from the resistances of the two sides of the circuit and the power supply current I p by the equations ##EQU2## which can be rearranged as ##EQU3## The switch resistance which can be generated by heating a length of tape superconductor is governed by the resistivity of the copper stabilizer in the tape, which has a resistivity ratio of about 75 which indicates a low temperature (4-20K) resistivity of about 2.27e-8Ω/cm. The copper, preferably, is 0.015 cm thick and 0.3 cm wide, so the resistance per unit length is 5μΩ/cm. Therefore, a conventional 1 inch heater will produce a resistance of about 13μΩ. Note that the tape temperature is below the transition temperature very close to the boundary of the heater because of the excellent thermal conductivity of the copper in the tape, so the length of the normal region in the switch is essentially equal to the length of the heater. Since the coil resistance is unknown, the fraction of the current which circulates in the coil will also be unknown. However, a desirable coil resistance is in the nΩ range. Using 13 nΩ as a working value, the coil current will be 0.999I p . Should the resistance of the coil be significantly higher then 13nΩ, it will not affect the ability of the apparatus to yield an accurate resistance measurement. This is determined solely from field drift rate and inductance. But the current at which the resistance is measured will be a smaller fraction of the power supply current. Even at a coil resistance of 1.3 μΩ, the coil current will still be about 0.91I p . This is not seen as a significant drawback to this testing method. The placement of the field sensor is an important part of the proposed testing technique, because there are fields generated by means other than the transport current in the superconductor. These means include circulating currents which are driven by changes in magnetic flux in the region of the conductor and Meissner effect currents. Fortunately, these fields drop off as a higher power of the distance from the coil than does that produced by the transport current. In order to ensure that the effects of the fields from circulating currents are minimal, the field sensor must be placed sufficiently far from the coil that the transport current field dominates the reading. For instance, while the center of the switch loop is an attractive place to position the sensor, the field from circulating currents in the main coil at that position may be substantial. Since the bifilar coil actually produces a finite field at a distance away from the coil where the circulating current field is small, the field of the loop is not required for the measurement. In fact, a location distinct from the center of the loop is preferred. Thus far the geometry has been presented in which the loop extension from the coil is in a plane perpendicular to the axis of the bifilar coil. While this is the easiest winding geometry, especially for tape conductor, it presents a difficulty because the principal fields of the bifilar coil and the loop extension are in the same direction. With respect to FIG. 4, in order to minimize the effects of the field of the bifilar coil, the loop extension or switch 28 may be aligned with its plane perpendicular to the direction of arrow C along the azimuthal (circumferential) direction (the direction of arrow B of the bifilar coil winding 21). Since the bifilar coil field has only radial and axial components, Hall sensor 38 measures only the field of the loop extension or switch 28, which is preferred. Three other preferences for the loop extension 28 location relative to the coil 21 are: 1. radially far from the coil winding 21 (provided a large radius dewar is available); 2. axially far from the coil winding 21 (provided a large height dewar is available); 3. reaching to a location which allows placement of Hall sensor 38 on the bifilar coil winding axis. With respect to FIG. 5, the measured field contributions from Meissner effect circulating currents in the tape superconductor comprising the loop extension or switch 28 can be eliminated by coupling a ring 52, preferably, made of ferrite, mild steel, or other suitable ferromagnetic material to the loop extension or switch 28. Ring 52 has a gap 56 in which a Hall sensor 38 could be placed to measure the flux circulating in the ring, without being affected by the field of the bifilar coil winding 21 or the circulating currents in the loop extension or switch 28. The current sensitivity expression for the loop extension is: ##EQU4## where I is the transport current in the loop, dl the differential length along the iron ring, B the flux density in the ring, μ rel the relative permeability of the ring, and μ o is the magnetic permeability of free space. The calibration of ring 52 is achieved by winding a coil 54 on the ring 52 in a toroidal fashion such that any magnetic flux change in ring 52 will induce a current in the toroidal coil 54. Coil 54, preferably, is constructed of copper. Note that the toroidal or calibration coil 54 need not be a complete toroid, but may be only cover part of the circumference of the ring 52. In the linear range of a conventional B-H curve for the ferrite or soft iron ring 52, the toroidal current is linearly dependent on the current in the loop extension 28. The same toroidal calibration coil 54 can be used to null out the field of the loop extension or switch 28 for the purpose of improving the accuracy of measuring the deviation from that null. In this configuration, the temporal stability of Hall sensor 38 can be enhanced by reducing the contribution of loop extension or switch 28 and by using a conventional highly stable calibrated current source, such as, is used in the powering Hall sensor 38 itself. A temporal stability of about 10 ppm has been achieved in this manner. With respect to FIGS. 6a-6c, an improvement in the measurement and the information which may be derived therefrom can be realized by operating all or part of the bifilar coil winding 21 in a background magnetic field 104 of background field assembly 100. Since the conductor under test will eventually be wound into a field-producing coil, it will operate in a field. The characteristics of a superconductor vary depending on the field, so it is best to measure the resistance of the test length of conductor in a background field 104 which approximates that in which it will operate. Several alternatives for producing such a field with a solenoidal coil 102 are described here: 1. Place solenoid 102, which is small relative to the bifilar coil, winding 21 against the side flange of the bifilar coil winding 21, thus exposing part of the bifilar coil winding 21 to field 104 (FIG. 6a); 2. Insert the bifilar coil winding 21 in the bore of a large solenoid 102 (FIG. 6b); 3. Insert the bifilar coil winding 21 in the radial gap between two solenoids 102a and 102b which produce fields 104a and 104b, respectively, in the same direction in the gap (FIG. 6c). This option allows a fairly low field 104a and 104b to be generated in the central region of the interior field coil, where very convenient measurements free of background field can be made that include all features of the Hall sensor 38, loop extensions or switch 28 and rings 52. The stabilizing of the background field coil current can be accomplished by using a commercially available highly stable power supply (available to stability levels of 10 ppm or better) or a superconducting switch. In fact, an acceptable level of insensitivity to field fluctuations can be achieved with a standard power supply (100 ppm or so) because the mutual inductance between the background field solenoid 102 and the bifilar coil 21 is reduced by the ratio of the volume of the bifilar winding to the cylindrical volume enclosed by the bifilar coil 21. This ratio is about 4%, so a 100 ppm stability in the power supply for the background field coil 102 translates into a 4 ppm stability of the bifilar persistent current loop. This value is much better than the resolution of the Hall sensor measuring device and the ferrite coil coupling. Once given the above disclosure, many other features, modifications and improvements will become apparent to the skilled artisan. Such features, modifications and improvements are, therefore, considered to be part of this invention, the scope of which is to be determined by the following claims.

This invention relates to an apparatus and method for measuring the resistance of superconductors. Structures of this type, generally, allow the resistance of the superconductor to be accurately measured in a non-destructive manner by using a bifilar coil which includes an integrated loop/switch formed from the bifilar coil.

TECHNICAL FIELD The invention relates to artificial intelligence (AI) production systems, and more particularly, to a method for optimizing the pattern-matching phase of a cyclic, rule-based, data-sensitive AI production system. BACKGROUND It should be appreciated that the artificial intelligence branch of computer science has exhibited explosive growth in recent years. One facet of AI has been concerned with the modeling and use of inference systems. Such inference systems exploit the computer science formalism termed "rewrite" or "production" systems. At a minimum, a rewrite or production system includes an alphabet, an initial set of strings (axioms) defined over the alphabet, and a set of production or rewrite rules from which new strings (theorems) can be obtained. Peter Jackson, "Introduction to Expert Systems", Addison-Wesley Publishing Co., copyright 1986, pp. 29-51 and 126-141, points out that an AI production system comprises a rule set (sometimes called production memory), a rule interpreter that decides how and when to apply the rules, and a working memory that holds data, goals, and intermediate results. Brownston et al, "Programming Expert Systems in OPS5", Addison-Wesley Publishing Co., copyright 1985, pp. 4-31, contrast a rule-based production computational model with that of a procedural model. More particularly, Brownston describes a production system as an executable set of production rules in which each rule in such a system represents an ordered pair of a condition (state set) and action statements. This leads to Brownston's characterization of a production system as an unordered finite sequence of data-sensitive production rules. ______________________________________Brownston's ComparisonProduction System Model Procedural Modelof Computation of Computation______________________________________Program Description: Program Description:A description of data An ordered list ofexpressed as objects, instructions written inattributes, and values, a language with a well-and an unordered finite defined syntax andsequence of rules that semantics. The list hascan be referenced by the a specified beginning.data. Each rule consists The language includes aof a condition/pattern stop/halt instruction orpart and an action part. punctuation whose meaning is to cease instruction execution.Execution: Execution:Requires maintenance of a The instructions areglobal data base containing directly executable.the problem description and The first instructionany modifications, in the list initiatesadditions, or deletions execution. After this,thereto and a recognize, execution proceeds inresolve, act (RRA) or match, sequence punctuated byselect, execute cycle. The conditional branchescycle: until a stop or halt instruction is encountered.(a) identifies that subsetof rules having acondition or pattern partmatching the data,(b) selects at least onerule from the identifiedsubset of rules accordingto an extrinsic protocol,and(c) executes (fires) theaction prescribed by theaction part of theselected rule includingmodification to the database.______________________________________ In addition to Jackson and Brownston, reference should also be made to: (1) Miranker, "TREAT: A Better Match Algorithm for AI Production Systems", Proceedings of the AAAI-87 Sixth National Conference on Artificial Intelligence, Vol. 1, July 13-17, 1987, pp. 42-47. (2) Miranker, Dept. of Computer Science, University of Texas at Austin, Report TR-87-03, January 1987. (3) Forgey, "OPS5 Users Manual", CMU-CS-81-135, copyright 1981. (4) Forgey, "Rete: A Fast Algorithm for the Many Pattern/Many Object Pattern Match Problem", Artificial Intelligence, Vol. 19, copyright 1982, pp. 17-37. (5) Schor et al, "Advances in Rete Pattern Matching", Proceedings of AAAI '86. (6) Chambers et al, "Distributed Computing", Academic Press, copyright 1984, pp. 10-19. (7) Aho et al, "Compilers: Principles, Techniques, and Tools", Addison-Wesley Publishing Co., copyright 1986, pp. 608-632. Brownston, Miranker, Forgey, and Schor describe pattern-driven, forward-chaining production systems using a matching-rule, select-rule, execute-rule cycle based on the OPS5 AI language. Furthermore, these references point out that the process of many data-object/many pattern matching is the most computationally intensive phase in the production system control cycle. The references teach several techniques for reducing the computational intensity of the pattern-matching phase. First, advantage can be taken of the temporal redundancy of the data objects resident in working memory. That is, the set of objects over which the pattern portions of the rules are compared can be limited to those objects in working memory which have been either created or modified since the last cycle. Second, the many object/many pattern comparison can be systematized through use of sorting or dataflow-like processing networks. While there are other comparison algorithms, such as Miranker's TREAT, the best known comparison method having these attributes is the RETE algorithm ascribed to Forgey. Forgey and Jackson discuss the RETE pattern/object matching method used in the AI production system control cycle. This method includes: (a) compiling the condition elements of the pattern portion of a rule into an augmented data flowgraph (see Aho and Chambers) or a comparison sorting network (Forgey); (b) comparing each object with conditions of the pattern as expressed in the compiled network over a set of nodes (alpha-nodes) that each test one object at a time; (c) passing tokens indicative of a match from antecedent nodes to descendant nodes (beta-nodes) joined on a pattern-determined basis, comparing each token received at a descendant node, and passing tokens on to descendants in turn until the paths through the flowgraph are traversed; and (d) maintaining a list of instantiations satisfying the match conditions expressed at each node. In the remainder of this specification, the terms MAKE and CREATE, MODIFY and UPDATE, and REMOVE and DELETE will be used interchangeably. Also, the RECOGNIZE, RESOLVE, and ACT phases of the AI production system control cycle serve as synonyms for MATCH, SELECT, and EXECUTE. SUMMARY OF THE INVENTION It is an object of this invention to devise a method for optimizing the pattern-matching phase of a cyclic, rule-based, data-sensitive production system. It is a related object to devise a method in which a path through a RETE network used during the comparison phase can be traced, thereby minimally invoking the comparison, token passing, and recording aspects of the RETE algorithm. It is yet another object to devise a method utilizing the path for deletion of an object, including removal of any counterpart instantiations recorded at the lists maintained by various nodes through which a token has passed. It is a further object to devise a method in which alteration of objects by a MODIFY command is accomplished as a function of the DELETE and MAKE commands and their associated processes. The aforementioned objects are satisfied by the method steps executed during the pattern-matching portion of the matching, selection, and execute cycle of an AI production system comprising: (a) compiling a RETE network of the condition elements of the pattern portion of the rule being matched, the join nodes of said network being grouped in a pattern-determined associative manner; and (b) applying those data objects created or modified in the immediately preceding cycle to said RETE network. The method includes the further steps at each node of: (b1) maintaining a list of instantiations satisfying the match conditions expressed at that node, (b2) passing tokens to descendant nodes upon an object/pattern comparison match, (b3) maintaining pointers to all ancestor nodes through which the token for each object passed, and (b4) traversing said pointers as a path for avoiding those RETE node pattern/object matchings redundant between a previously matched object and an object being processed. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 depicts a directed graph RETE network produced during compilation of the pattern portion (left-hand side) of the sample rule. FIG. 2 shows a block diagram of the prior art logical machine as executed on any general purpose, stored program-controlled digital computer upon which the method of the invention may be practiced. DESCRIPTION OF THE PREFERRED EMBODIMENT In order to enhance appreciation for the method of this invention, a description of the generation, function and use of the RETE comparison algorithm and supporting data structures according to the prior art will be set out. Following this is a discussion of the invention utilizing the data structures. Lastly, a pseudo-code implementation and description is presented. Classic RETE Processing Reference should be made to previously cited Forgey, Brownston, and Miranker publications for implementation and additional details. The example discussed is a scheduling application involving parts and machines. Suppose there is a company that manufactures parts, that the company has two machines A and B, that there are two types of parts (P and Q) manufactured by the company, and that to complete a part of type P, machine A and then machine B must be used; while to complete a part of type Q, machine B and then machine A must be used. Suppose there is interest in a computer program that does production scheduling (deciding which parts should be routed to which machines) as the work progresses. Machines will be represented by records matching the following data declaration: ______________________________________ Declare machine, type, state.______________________________________ Here "type" and "state" are attributes or fields in the machine record. "Type" may have the value "A" or the value "B"--representing the type of machine, and "state" may have the value "AVAILABLE" or "BUSY"--representing the current availability of the machine to accept a part to begin processing. In this example, data-driven production system techniques as well as terminology common to such systems are used. Specifically, many records will be able to co-exist, each satisfying the declaration above. The collection of all records satisfying any one such declaration will be called a "class". Each individual record will be called a "class member". The overall collection of all class members from all classes will be called "working memory". Each part (partially or fully complete) will be represented in the system by class members (records) satisfying the following declaration: ______________________________________ Declare part, type, state, next-machine-needed.______________________________________ Here "type" can have the values "P" or "Q", while "state" must be "in-process", "waiting", or "complete", representing that a part must at any instant be either in the process of being machined, or waiting for a machine, or through with all needed machining. Suppose also that occasionally the company gets express orders that must be expedited. Such orders are of very high priority, and although the machining of a nonexpress part will not be interrupted to make the machine available for an express part, the company's policy is that while there is an express part in the system that will need a machine in the future, that machine will not start any nonexpress work--thus increasing the likelihood that the machine will be available or near becoming available when it is needed for an express part. The declaration for a class member that would represent an express order might look like the following: ______________________________________ Declare express-order, part-type, state, array-of-all-machines-needed, index-of-next-machine-needed.______________________________________ Using data-driven, forward-chaining, condition-action rules (or productions) to encode the routing of parts to machines, one rule in the computer program might look like the following when paraphrased in English: ______________________________________Sample-rule:Whenthere is a machine - which shall be called Msuch that M is availableand there is a part - which shall be called Osuch that the part is waiting for its nextmachining step andthe next machining step required for O must bedone on the machine Mand there are no express parts that satisfy boththe part is not complete, i.e., it is waitingor in-process andthe part will need machine M before it iscompleteThen take the following actions:change the state of M to BUSYchange the state of O to IN-PROCESSroute the part O to the machine M for processing.______________________________________ The list of conditions between the `When` and the `Then` in the rule is called the antecedent or left-hand side (LHS) of the rule. The actions to be taken when the LHS is satisfied by some list of class members is called the consequent or right-hand side (RHS). There would be many other rules in such a system, and there would be a method for notifying the computer program whenever a machine completed its work on a part. Referring now to FIG. 1, there is shown a RETE network for the LHS of the Sample-Rule. All arcs are directed and are considered as pointing down in the above picture. RETE processing starts at the top nodes (class anchors) and flows downward through the network. The nodes marked MACHINE, PART, and EXPRESS-ORDER are, respectively, anchors for queues of all machine, part, and express-order class members that exist in the program at that point in time. Changes to machine, part, or express-order class members can be of several kinds. New members can be created, representing new objects that the system should consider in routing parts to machines. Class members can be destroyed, representing the deletion of some object from participation in the system. For example, when a part is shipped to a purchaser, it might be natural to delete the class member representing that part from the system, as it should no longer be considered in any decisions about scheduling. Changes can also be made to existing class members. Different values are assigned to the attributes of the class member. This is done to represent changes in state, etc. Each time a change is made to working memory that change is passed through the RETE network. In rough terms, the central nodes of the network fragment in FIG. 1 correspond to tests that must be performed to determine what lists of class members satisfy all the conditions in the LHS of the Sample-Rule above. Class members arranged in lists are tested at a node to determine which lists satisfy the corresponding conditions and, for each node, summary information is retained in a queue based at that node. Usually, the retained information indicates which lists of class members passed the test associated with the node. The nodes labeled AVAILABLE?, WAITING?, and NOT-COMPLETE? are called alpha-nodes. Alpha-nodes correspond to tests that only mention a single class member. These nodes have only one incoming arc in the RETE network. The nodes labeled PART-NEEDS-MACHINE? and MACHINE-RESERVED-FOR-EXPRESS? are called beta-nodes. In this RETE network, the beta-nodes all have two incoming arcs and correspond to tests that determine whether a list of several class members satisfies a condition. For example, the node labeled PART-NEEDS-MACHINE? corresponds to a test that determines whether a pair of class members, the first a part and the second a machine, is such that the next-machine-needed field on the part class member is identical to the type filed of the machine class member. In other words, when given a part and a machine, the node determines whether the part needs to go to that machine next. If a pair (or, in general, a list of any length) of class members satisfies a test, that fact is recorded in a control block which shall be called a Satisfaction Block (SB). The SB is placed in a queue that is anchored at the RETE node and becomes a record of what lists of objects satisfied the conditions associated with that node. The declaration for a simple SB might look like the following: ______________________________________ Declare satisfaction-block, next-SB-in-queue, list-of-class-members;______________________________________ where the specified list-of-class-members field indicates a list of class members which together satisfy the test associated with the node. Lists of objects that satisfy the conditions associated with a node are candidates that need to be considered for satisfaction of later conditions. To make all the needed tests, while restricting consideration to just those lists of class members that have satisfied all earlier tests, a computational strategy is adopted that pushes tokens through the RETE network (top to bottom when the network is drawn, as in FIG. 1). A token represents a change to some SB and reflects processing that must be done at a successor node. Tokens might satisfy the following declaration where the phrases following the `/*` on any line are comments: ______________________________________Declare token,change-type, /* MAKE or DELETElist-of-class-members, /* list of one or more class /* membersnode, /* identifies node where /* tests will be attempteddirection; /* ancestor that passed the /* token is LEFT or RIGHT of /* the node in the token______________________________________ The declaration of a token includes a direction field which can be set to have one of two values, LEFT or RIGHT. It shall be assumed that all the RETE networks under discussion satisfy the following principles: A node with only one incoming arc is always viewed as having its one ancestor as a left ancestor. Where a node has two ancestors (a beta-node) and each ancestor is an alpha-node or is an anchor to a class, then one ancestor can be arbitrarily picked as the left ancestor, while the other will be the right. Where a beta-node has one ancestor that is a beta-node, then that ancestor shall be called the left ancestor, and the other ancestor (which in the example and in most implementations must be an anchor or alpha-node) becomes the right ancestor. The use and flow of tokens are illustrated by examining the processing that would take place in the sample RETE network of FIG. 1. Suppose the process is started with an empty working memory. The RETE network of FIG. 1 does not contain any satisfaction blocks. Further, suppose that during initialization of the execution of the program containing the rule above, the first machine class member is created, representing machine A. Since there is no work yet for the machine to do, the state field on the class member would be set to AVAILABLE. For each successor in the RETE network of the machine anchor node, a token would be created that contains a list of one class member--the new machine class member--in its list-of-class-members field. In the current example, there is only one node that is a successor of the machine anchor node. Thus, only one token would be generated, and it would point to the new machine for its list-of-class-members field. It would point to the AVAILABLE? node for its node field, and it would specify that the token arrived at the AVAILABLE? node from a LEFT ancestor. The token would contain MAKE in its CHANGE-TYPE field, indicating that the new class member was just created. Typically, if there are many tokens generated at once, they are placed on a stack and processed one at a time. In this simple case, only one token has been generated, so the processing for that one token would immediately proceed. (Typically, the token would be pushed on the token stack but soon thereafter would be popped off.) This processing done at the node would include executing the test associated with the node for the machine (class member) in the token. In this case, the newly created machine would be tested to see whether its state field listed it as available--which it would be. Thus, the test associated with the node would be passed and several actions would be taken. First, a Satisfaction Block would be created to record that the list of class members in the token passed the test. In this case, the new SB would point to the one new machine class member. Second, new tokens with change-type of MAKE would be generated, one for each successor of the AVAILABLE? node. Again, in this simple case, there is only one successor to the AVAILABLE? node so a token is stacked that specifies (a) machine A, (b) the PART-NEEDS-MACHINE? node, and (c) LEFT. Again, since there is only one node on the stack, processing for that node immediately proceeds. However, a list of a machine and a part is required to satisfy the PART-NEEDS-MACHINE? node, and there are no parts in existence yet. Thus, there is no more processing that can be done at this time, so control is returned back to the application. Suppose that the next action taken by the application during this initialization phase is to create the second machine; namely, machine B. As before, a token would be generated to cause testing at the AVAILABLE? node to determine whether this machine is available--which it would be. An SB is created so that there are now two SBs in the queue off the AVAILABLE? node. Since nothing more can be done, control (as before) returns to the application. Suppose that the application begins normal operation, and scheduling for the first part is to start. The application makes a part class member to represent the existence of the first part, a type P part that is waiting to be machined. Since type P parts must first be machined by machine A, the part's NEEDS-MACHINE field is set to A. Again, a token with change-type of MAKE is generated, indicating this change. The token would point to the new part, specify the WAITING? node, and the left direction. Processing of the token at the WAITING? node would discover that the part was indeed waiting, so an SB would be created and enqueued off the node and a new token would be created (and pushed onto the stack), indicating that this new part needs to be tested at node PART-NEEDS-MACHINE? where the arrival is from the right. With only one token on the stack, it is popped off and processing for that token proceeds immediately. A machine and a part are now needed to perform the test associated with the condition at a beta-node (the PART-NEEDS-MACHINE? node). There are two machines, A and B, which have previously arrived in tokens at the PART-NEEDS-MACHINE? node. This fact is recorded by the existence of two SBs enqueued off the AVAILABLE? node, which is the left ancestor of the current node. Thus, a loop is executed that walks through all SB blocks off the AVAILABLE? node. For each such SB, the test is performed at the PART-NEEDS-MACHINE? node to determine if the new part (of type P) needs the machine in the SB for its next machining step. If the test (PART-NEEDS-MACHINE?) fails, then no additional action is taken. If the test passes, then a new SB is created and enqueued off the PART-NEEDS-MACHINE? node. In this case, the test of the type P part with machine B fails, and the test with machine A passes. Thus, an SB with a list of class members comprised of machine A and the new type P part is enqueued off the PART-NEEDS-MACHINE? node, and a token representing this change to PART-NEEDS-MACHINES?'s SB queue is generated for the one successor of this node. The change made to PART-NEEDS-MACHINE?'s SB queue is the creation of a new SB. Thus, the new token is a MAKE token, and its list of class members includes the A machine and the new P part. PART-NEEDS-MACHINE? is an example of a positive beta-node. When a MAKE token arrives at such a node from one direction, the list of class members in the token is successively augmented with the list of class members from each SB off the predecessor node in the opposite direction. For each such SB, the test for the current node is performed using the augmented class member list. A passed test causes more tokens to be generated--one for each successor node. These new tokens are MAKE tokens that represent new candidates at successor nodes. This token is popped off the stack, having arrived at the MACHINE-RESERVED-FOR-EXPRESS? node from the left. This node is a negative beta-node. Unlike positive beta-nodes, the output of a negative beta-node is not an augmented token. A token arriving at a negative beta-node from the left includes a list of class members. This same list is either passed intact to all successor nodes, or it is stopped altogether. The list is passed on (in tokens) if no test passes at this node involving this class member list augmented by any one of the lists in the right ancestor's SB queue. This list is stopped and not passed to any successor if some test passes at this node when the arriving list is augmented in turn by each of the lists from the right ancestor's SB queue. Negative beta-nodes are often discussed in the following terms. Class member lists that arrive from the right will stop class member lists arriving from the left from passing through. Thus, a negative beta-node is much like a gate, where the things arriving from the right (according to what tests pass and fail) determine whether things arriving from the left will pass through. It should be noted that the arrival from the right of a DELETE token (indicating the deletion of an SB from the right ancestor's queue) at a negative beta-node can cause creation of a new SB at the negative beta-node and the corresponding generation of MAKE tokens for all successors of the negative beta-node. This happens when the arriving token indicates deletion of the only SB in the right ancestor's queue that was stopping some SB in the left ancestor's queue. Likewise, a MAKE token arriving from the right can stop an SB that is in the negative beta-node's SB queue, and the MAKE token can cause the removal of an SB and the generation of DELETE tokens for all successors. In this example, the right ancestor of the MACHINE-RESERVED-FOR-EXPRESS? node is the NOT-COMPLETE? node. There are no express orders, so there can be no express orders with the part not complete. Therefore, there is nothing in NOT-COMPLETE?'s SB queue, and there is nothing that will stop any list of class members (in a token) arriving from the left from passing on through. Thus, arrival of a MAKE token from the left causes a new MAKE token to be generated and passed to the production node. There is only one production node shown. However, in a real RETE network, there is one production node for each rule. Whenever a token arrives at a production node, the arrival indicates some change to the conflict set for that rule. A MAKE token arriving at a production node indicates that a new instantiation should be created. A DELETE token arriving indicates that an existing instantiation should be eliminated. In this example, the arrival of the MAKE token at the production node indicates that a new instantiation should be made. Indeed, the pair of class members consisting of machine A and the one part of type P satisfies the conditions of the Sample-Rule based on all the class members that exist at this time. Thus, this one instantiation would normally become a candidate for firing with all other instantiations of other rules. If the rule did fire, then the action part of the rule would change the status of both the P type part and of machine A. Those changes, when pushed through the RETE network, would invalidate this on instantiation to Sample-Rule. Next, the RETE processing is examined in more detail for this case. As expressed before, this example merely illustrates the prior art. Suppose that the single instantiation is selected and Sample-Rule does fire. The first action taken is to change the status of machine A to BUSY. This change in machine A's state field with the classic RETE algorithm would be treated as a deletion of the old machine class member, followed by the creation of a new machine class member that is identical to the old one except for the altered value in the state field--the new value being BUSY. If the machine A class member is deleted, then before the class member is destroyed, a DELETE token is created for the AVAILABLE? node with the token pointing to the machine A class member. As always, the token is stacked with any other tokens. When popped and processed for the AVAILABLE? node, the test associated with the AVAILABLE? node is repeated to determine whether the class member passed or failed the test when it was created. If the current test passes (or fails), then the original test must have likewise passed (or failed, respectively) when the MAKE token was processed. If the earlier test did fail, there was no additional RETE processing generated by the token. No SB block was made, and no new tokens were spawned. Thus, there is nothing more that need be done now, as there are no descendant references in the RETE network that must be removed to reflect deletion of the class member. If the earlier test passed, then there must be a record of the fact in the SB queue for the node. This is found by searching the queue, i.e., walking through all the SB blocks in the queue until one is found with a class member list that is identical to the class member list for the current token. This SB block is then dequeued and destroyed. Also, the fact that the test passed indicates that additional tokens were made when the machine A class member was created, and those tokens must be sought out and destroyed since they mention (typically, they point to) a class member that is being destroyed. A new token (a DELETE token) is created for each successor to the current node, and these are all placed on the stack for later processing. In summary, processing for a DELETE token exactly token. All tests are repeated. When a test is passed, there is an additional expense of searching for the right SB to excise. When a DELETE token arrives at a positive beta-node from one direction, it must be paired with all class member lists from SBs enqueued off the opposite ancestor. For each test that passes, the SB in the current block must be found and excised, and new tokens must be sent to successors. Processing of a DELETE token at a negative beta-node likewise undoes the work of a MAKE token. The computational expense of processing DELETE tokens is greater than that for processing MAKE tokens. It is frequently necessary to do an additional search of the SB queue at a node to find and excise the appropriate SBs. RETE Processing as Modified According to the Invention It was unexpectedly observed that if additional information links were maintained in an SB with its successor SBs, then those SBs that mention the same class member list as in the original SB could immediately be deleted if the original SB were deleted. The process for creating new SBs is largely unchanged, other than the responsibility for saving the additional information. However, in the resulting structure, a class member is the root of a tree of all SBs (and instantiations) that mention the class member, the class member itself being treated like an SB enqueued off an anchor node. Thus, deletion of all SBs (and instantiations) that mention the class member is largely a matter of walking a tree and excising all SBs encountered. The algorithm can be implemented so that no test is ever repeated. This reduction in the total number of tests to execute RETE processing is especially important when tests are slow lo execute; for example, when class members are stored on a slower storage medium, and accessing class members for testing is computationally expensive. Also, with the invention described herein, no searching of SB queues is required to find the SBs to be excised. This can markedly speed up execution of RETE processing, depending on the length of the SB queues. In this invention, an SB might have the following declaration: ______________________________________Declare improved-satisfaction-block,next-sat-block-in-queue /* used to enqueue this /* block in the doubly-previous-sat-block-in-queue /* linked list of all /* SBs for this nodelist-of-class-members, /* list of members /* passing test of nodeleft-descendant-SB, /* points to any left /* descendant or is /* nullleft-ancestor-SB, /* left ancestor /* pointer, null if no /* lf ancnext-left-sibling-SB, /* doubly-linked list /* of all SBs with theprevious-left-sibling-SB, /* same left ancestorright-descendant-SB, /* points to any right /* descendant, or nullright-ancestor-SB, /* right ancestor pntr, /* null if no rght ancnext-right-sibling-SB, /* doubly-linked list /* of all SBs with theprevious-right-sibling-SB; /* same right ancestor______________________________________ If node N is a successor of node M in a RETE network (i.e., there is an arc from M to N), and if X is an SB queued off M, and Y is an SB queued off N so that the class member list in X is an initial or terminal segment of the list in Y--so that the deletion of X would cause (using classic RETE processing techniques) tokens to be generated that would result in the deletion of Y, then say that Y is an immediate descendant or a successor of X, and say that X is an immediate ancestor or a predecessor of Y. If M is a left (or right) predecessor of N, then say that X is a left (or right, respectively) predecessor of Y, and Y is a left (or right, respectively) successor of X. If N is a beta-node, there may be many SBs enqueued off N that are all immediate descendants of the same SB enqueued off M. However, a given SB can have at most one left and at most one right immediate ancestor. Fields similar to those added to the SB declaration in order to arrive at the improved SB declaration (but excluding ancestor pointers) can be associated with each class member. Thus, in all particulars, a class member itself can be treated like an SB insofar as handling creations and deletions. Likewise, fields similar to those in an improved SB (excluding those that point to descendants) can be added to instantiations, as instantiations record what lists of class members pass through production nodes and they are similar to and treated like SBs. The sibling fields in an improved satisfaction block are used to form a doubly-linked (usually circular) list of all SBs that have a common immediate ancestor. The left and right descendant fields point respectively to any left/right descendant. The left and right ancestor fields are each either set to point to the appropriate ancestor SB if such exists, or are set to a null value if no such ancestor SB exists. When a MAKE token is being processed at either an alpha or a positive beta-node and a test is passed, all ancestor SBs are known and all the fields in an improved SB can be easily set to the correct values. When a delete of a class member is being processed, a depth-first walk of the descendant SB tree can be executed, and each SB can be deleted from the doubly-linked queue of sibling SBs. If that deletion exhausts the sibling queue, then the predecessor SB's descendant field can be set to null. If the predecessor SB's descendant field points to the deleted SB and the sibling queue is not depleted by the deletion, then the ancestor SB's descendant field can be set to point to another descendant SB. Once an SB is no longer referred to (pointed to) by any other SB, it can safely be destroyed. Care must be taken to eliminate all left and all right descendants of any SB that is going to be deleted. Processing at negative beta-nodes can be done in several ways. However, it is not possible to simply follow an immediate translation of the above-described techniques into the world of negative beta-nodes. The problem in using the classic approach to processing at negative beta-nodes is that a complete record is not maintained of the result of every test performed at a negative beta-node. With the processing for positive beta-nodes, as described above, every passing test is recorded in an SB. A failing test is recorded by the nonexistence of an SB. In contrast, the SB queue for a negative beta-node (with the classic algorithm as earlier described) is a subset of the SB queue for the left ancestor node. An SB is created and kept (off the negative beta-node) exactly in those cases where no SB in the right ancestor's queue passes the test when paired with the SB from the left ancestor. Thus, the results of all tests involving SBs that appear in the negative beta-nodes's SB queue are known; namely, all such tests failed. However, if an SB from a left ancestor does not reappear in the negative beta-node's SB queue, then there is no information retained about which SB or SBs from the right ancestor node stopped the left arriving SB by passing the test when the left and right SBs were coupled together. Since negative beta-nodes tend to be significantly less common than positive beta-nodes, it is possible to treat them using an approach similar to that in the classic RETE processing while still obtaining marked speed improvements for most applications. Using this approach, the elimination of an SB from a node that is a right ancestor of a negative beta-node would trigger the spawning of DELETE tokens that would be passed on to the negative beta-node and processed in the classical way. It is also possible to keep additional information at each negative beta-node about which SBs from the right stop which SBs arriving from the left. This can be done, for instance, by keeping a list (the standard SB queue for the negative beta-node) of all left arriving SBs that are not stopped, and also by keeping another list of all left arriving SBs that are stopped and keeping with each entry in that list another list of all the right arriving SBs that stop the associated left arriving SBs. Blocks in this new list might be called Stopping SB Blocks, or SSBBs. In this case, fields similar to those in the improved-satisfaction-blocks (outlined above) can be used to link SSBBs to their ancestor SBs. While the details of the algorithm are different for processing at positive and negative beta-nodes, the general approach is similar to that outlined above for the processing done a alpha and positive beta-nodes. Pseudo-code Implementation and Comment One use of the RETE algorithm and of the method of this invention is in the implementation of the pattern matching for data-driven production systems. The control cycle for such a system is shown with reference to FIG. 2. In that environment, when a program is first initialized, the RETE network for the set of rules in the application must be built. However, the invention relates to modifications to a RETE network once built, and the following discussion assumes that one has already been compiled. The fundamental loop or cycle of execution for a data-driven production system is the Match-Resolve-Act cycle. This is illustrated in FIG. 2 and set out in the pseudo-code outline as follows: ______________________________________Some initial conflict set (set of all instantiations)is given;Do Forever;Call conflict-resolution;/* select the best instantiation to fireIf there is a best instantiationThen Do;Fire the best instantiation;/* Note that during the course of firing the/* instantiation, the action code will initiate/* changes to working memory. The application/* code will call MAKE, DELETE, and MODIFY to/* notify the RETE algorithm about changes to/* working memory, and the RETE algorithm will/* push these changes through the RETE network/* in order to compute the resulting changes/* to the conflict set.MAKE, DELETE, and MODIFY may be called by theaction part of a rule during the execution ofthe action part of the rule;End;Else Do; /* the conflict set is emptyReturn to the user of the application and allowhim to terminate execution or to interactivelychange working memory and therein initiate callsto MAKE, DELETE, or MODIFY;End;End;______________________________________ The heart of the RETE processing is thus in the routines MAKE, DELETE, and MODIFY. These routines are notified of changes to working memory, and they push those changes through the RETE network and thereby compute the corresponding changes to the conflict set. An outline of these routines follows. As derived from the illustrative example in the previous section, the data structures that correspond to the following declarations shall be used: ______________________________________Declare tokenlist-of-class-members, /* list of one or more /* class membersnode, /* identifies node where /* tests will be attempteddirection; /* ancestor that passed the /* token is LEFT or RIGHT /* of the node in the token______________________________________ It should be noted that tokens will only be used for MAKE-type actions in the outlined implementation; thus, there is no change-type field in this declaration of a token. The declaration of a Satisfaction Block (SB) should be formatted as: ______________________________________Declare improved-satisfaction-block,next-sat-block-in-queue /* used to enqueue this /* block in the doubly-previous-sat-block-in-queue /* linked list of all /* SBs for this nodelist-of-class-members, /* list of members /* passing test of nodeleft-descendant-SB, /* points to any left /* descendant or is /* nullleft-ancestor-SB, /* left ancestor /* pointer, null if no /* lf ancnext-left-sibling-SB, /* doubly-linked list /* of all SBs with theprevious-left-sibling-SB, /* same left ancestorright-descendant-SB, /* points to any right /* descendant, or nullright-ancestor-SB, /* right ancestor pntr, /* null if no rght ancnext-right-sibling-SB, /* doubly-linked list /* of all SBs with theprevious-right-sibling-SB; /* same right ancestor______________________________________ The declaration of a Stopped SB Block (SSBB) consists of: ______________________________________Declare stopped-SB-block,next-SBBB-in-queue, /* used to enqueue this /* block in the doubly-previous-SSBB-in-queue, /* linked list of all /* SSBBs for this nodeleft-ancestor, /* points to stopped SBnext-left-sibling, /* doubly-linked list of /* all SBs and SSBBsprevious-left-sibling, /* with same left ancestorright-ancestor, /* points to stopping SBnext-right-sibling, /* doubly-linked list of /* all SBs and SSBBsprevious-right-sibling; /* with same right ancestor______________________________________ MAKE is passed a pointer to a just created class member and an identifier of the class. It computes the corresponding changes to the conflict set. ______________________________________MAKE(class-member,class-identifier);Locate the anchor for the class specified by theclass-identifier;Enqueue the newly created class member in the classmember queue for the class (the queue is based atthe anchor for the class);Create a token and push it onto the token stackwith class-member as its (length one) list ofclass memberswith the anchor for the class as its RETE nodewith LEFT direction;Call MAKE-PROCESS to do the RETE processing;/* the input to MAKE-PROCESS is the token stackReturn to the caller;End MAKE;______________________________________ DELETE is passed a pointer to the class member that is to be deleted. DELETE computes the corresponding changes to the conflict set. Thereafter, it frees up the storage being used by the class member. ______________________________________DELETE(class-member):Call DELETE-PROCESS(class-member); /* do the RETE /* processingFree class-member's storage block;End DELETE;______________________________________ MODIFY is passed a pointer to the class member that is to be changed. It is also passed a description of the field that is to change, and it is passed the new value that is to be assigned to the field. MODIFY will delete the old copy of the class member insofar as the RETE algorithm is concerned. Then it makes the indicated change to the class member, and does the RETE processing to reflect the creation of a new class member replacing the old one, but with the one field changed. ______________________________________MODIFY(class-member,class-identifier,field-identifier,new-value-for-field):Call DELETE-PROCESS(class-member);Assign new value to field to be changed inclass-member;Call MAKE(class-member,class-identifier);End MODIFY;______________________________________ Thus, it can be seen that the heart of the processing takes place in the two routines MAKE-PROCESS and DELETE-PROCESS. These routines are outlined in pseudo-code. First, consider the MAKE-PROCESS routine. As mentioned before, the input to MAKE-PROCESS is the token stack, and MAKE-PROCESS is driven by the token stack. When that stack is empty, then MAKE-PROCESS's work is complete and it returns to the caller. ______________________________________MAKE-PROCESS:Do while token stack is not empty;Pop the top token off the token stack and call itthis-token;If node of this-token is a class anchor (so tokenlist is a list of one SB and it is a classmember)Then Do;/* use SET-SB-AND-SEND-ON to create SB/* recording passed test and to generate/* tokens for each successor of anchorCall SET-SB-AND-SEND-ON(anchor node, list ofone class member - same as list in this-token, null, null);End;ElseIf node of the token is an alpha-nodeThen Do;Execute the test associated with the node forthe one class member in the class memberlist of this-token;If the test passedThen Do;/* use SET-SB-AND-SEND-ON to create SB/* recording passed test and to generate/* tokens for each successor of anchorCall SET-SB-AND-SEND-ON(node of this-token,class member list of this-token, leftancestor SB, null);End;/* else the test failed, and we do nothing/* more for this-tokenEnd;ElseIf node of the token is a positive beta-nodeThen Do;Do for each SB in the SB queue of theancestor in the opposite direction fromthat of the arriving token;Form an augmented list by concatenating thelist of class members from the SB withthe list of class members from this-token;Execute the test associated with this nodefor the augmented list;If the test passedThen Call SET-SB-AND-SEND-ON(node of this-token, augmented list, left ancestor SB,right ancestor SB);End;End;ElseIf node of the token is a negative beta-nodeThen Do;If token is from leftThen Do;Do for each SB in the SB queue of the rightancestor;Form an augmented list by concatenatingthe list of class members (one classmember in this case) from the SB withthe list of class members from this-token;Execute the test associated with thisnode for the augmented list;If the test passedThen create and enqueue an SSBB pointingto the left and right ancestor SBs;End;If no test passed among tests for all rightancestor's SBsThen Call SET-SB-AND-SEND-ON(node of this-token, list from this-token, leftancestor SB, null);End;Else Do; /* token is from the rightDo for each SB in the SB queue of the leftancestor;Form an augmented list by concatenatingthe list of class members from the SBwith the list of class members (oflength one) from this-token;Execute the test associated with thisnode for the augmented list;If the test passedThen Do;Create and enqueue an SSBB pointing tothe left and right ancestor SBs;If there exist no other SSBBs enqueuedoff this node that refer to the sameleft ancestor SBThen Do;Locate the SB in the queue off thisnode that has the same class memberlist as the SB from the leftancestor node;Call DELETE-PROCESS(pointer to the SBenqueued off this node - as locatedabove);End;End;End;End;End;ElseIf node of the token is a positive merge-nodeThen Do;Do for each SB in the SB queue of theancestor in the opposite direction fromthat of the arriving token;Form an augmented list by concatenating thelist of class members from the SB withthe list of class members from this-token;Call SET-SB-AND-SEND-ON(node of this-token,augmented list, left ancestor SB, rightancestor SB);End;End;ElseIf node of the token is a negative merge-nodeThen Do;If token is from leftThen Do;Do for each SB in the SB queue of the rightancestor;Create and enqueue an SSBB pointing tothe left and right ancestor SBs;End;If the right ancestor node's SB queue isemptyThen Call SET-SB-AND-SEND-ON(node of this-token, list from this-token, leftancestor SB, null);End;Else Do; /* token is from the rightDo for each SB in the SB queue of the leftancestor;Create and enqueue an SSBB pointing tothe left and right ancestor SBs;If there exist no other SSBBs enqueuedoff this node that refer to the sameleft ancestor SBThen Do;Locate the SB in the queue off thisnode that has the same class memberlist as the SB from the left ancestornode;Call DELETE-PROCESS(pointer to the SBenqueued off this node - as locatedabove);End;End;End;End;Else /* the node type must be a production node -/* all other possibilities have been/* exhaustedDo;Create a new instantiation made up of list ofclass members from the token and the ruleassociated with the node in the token;End;End of do while loop;Return to callerEnd MAKE-PROCESS;______________________________________ There are several references in the above to the utility routine SET-SB-AND-SEND-ON. This routine accepts a node, a list of pointers to class members, a pointer to a left ancestor (possibly null), and a pointer to a right ancestor (possibly null). It does two things. First, it creates an SB from the list of pointers to class members, and it enqueues that SB off the passed node. Second, it loops through each successor node in the RETE network of the passed node, and for each it creates and pushes a token. A pseudo-code description follows: ______________________________________SET-SB-AND-SEND-ON(parent-node,list-of-class-members,left-ancestor-SB,right-ancestor-SB):Create an SB with list-of-class-members and enqueueit off parent-node;Do for each RETE node successor of parent-node;Make a tokenwith node being the successorwith direction being the kind (left or right)of ancestor that the parent-node is to thesuccessorwith the list of class members beinglist-of-class-members;Push the token on the token stack;End;End SET-SB-AND-SEND-ON;______________________________________ It remains to outline the DELETE-PROCESS routine. This is called to delete all references in the RETE network to an SB, an SSBB, in instantiation, or a class member--which, for the purposes of the RETE algorithm, is a special kind of SB. DELETE-PROCESS is called and passed either a pointer to a class member or a pointer to an SB. If a class member is passed, then that class member is deleted from its class queue. If an SB is passed, then that SB is deleted from the SB queue in which it resided--while the integrity of all SB queues is maintained. Whether as SB or a class member is passed, DELETE-PROCESS proceeds to call itself recursively in order to eliminate all left and right descendant SBs. Thinking of all the descendants of an SB as forming a tree, the recursive invocation of DELETE-PROCESS actually walks that tree in a depth-first manner and just before leaving a node in that tree (which is an SB), that SB is excised from all queues in which it resides and the SB is destroyed (i.e., the block of storage is freed): ______________________________________DELETE-PROCESS(SB-to-delete):/* recursively invoked where SB-to-delete is either/* an SB or a class member to be deletedIf SB-to-delete has a left successor SBThen Call DELETE-PROCESS(that left successor SB);If SB-to-delete has a right successor SBThen Call DELETE-PROCESS(that right successor SB);Excise SB-to-delete from the doubly-linked queue ofall SBs off the node with which SB-to-delete isassociated;If SB-to-delete has a left sibling SBThen Do;If SB-to-delete has a left ancestorand if the left successor field of the leftancestor points to SB-to-deleteThen reset that left successor field to point toa sibling of SB-to-delete;Excise SB-to-delete from its doubly-linked queueof left siblings;Call DELETE-PROCESS(left sibling ofSB-to-delete);End;Else /* SB-to-delete is alone in the left sibling/*queueIf SB-to-delete has a left ancestorThen set the left successor field of that leftancestor to null;If SB-to-delete is an SSBB(and it is known that SB-to-delete has a leftancestor)and there are no other SSBBs enqueued atSB-to-delete's node that are left siblings ofSB-to-delete (i.e., same left anc)Then Do;Call SET-SB-AND-SEND-ON(current node, classmember list of SB-to-delete's left ancestor,SB-to-delete's left ancestor, SB-to-delete'sright ancestor);Call MAKE-PROCESS;End;If SB-to-delete has a right sibling SBThen Do;If SB-to-delete has a right ancestorand if the right successor field of the rightancestor points to SB-to-deleteThen reset that right successor field to pointto a sibling of SB-to-delete;Excise SB-to-delete from its doubly-linked queueof right siblings;Call DELETE-PROCESS(right sibling ofSB-to-delete);End;Else /* SB-to-delete is alone in the right sibling/* queueIf SB-to-delete has a right ancestorThen set the right successor field of that rightancestor to null;If SB-to-delete is an SB rather than a class memberThen free SB-to-delete's storage block (and sodestroy SB-to-delete);Return to caller;End DELETE-PROCESS;______________________________________ It will be further understood by those skilled in this art that various changes in form and detail may be made therein without departing from the spirit and scope of the invention.

A demand-driven AI production system utilizing a RETE network for comparison matching in a condition/data match, rule-selection, and rule-firing execution cycle in which the RETE network is modified to maintain a list of instantiations satisfying the match conditions expressed in each node of the RETE network, passing of tokens to descendant nodes upon a comparison match, maintaining patterns to all ancestor nodes through which the tokens have passed, and traversing the patterns as a path for avoiding those RETE pattern matchings redundant between a previous match and a current match in progress.

FIELD OF THE INVENTION The field of this invention relates to freezing plugs in pipe, and in particular, to freezing plugs in large and subsea pipe. BACKGROUND OF INVENTION Pipe freezing is an established technique for effecting the temporary isolation of a section of a pipe. Fluid in a line, which may include water, oil, hydrocarbons or gel and/or which may involve a strategically placed substance or pig, is frozen by surrounding the pipe with a cooling medium. Current industry practice follows the cryogenic fluid system for performing a pipe freeze. The cryogen typically comprises a liquid nitrogen or carbon dioxide. In the cryogenic system compressed fluids (or solids) at very low temperatures are supplied to a job in tanks, allowed to vaporize to absorb heat and are then vented to the atmosphere. For each job, a sufficient supply of the cryogen, in liquid or solid form, must be on hand to carry out the freeze. Certain cryogenic fluids, such as liquid nitrogen or carbon dioxide, readily meet the needs of vaporizing at very low temperatures and being generally cost effective and transportable. The discharge of these fluids, when consumed, is environmentally acceptable. The bulk of the tanks required to assure an adequate supply of the cryogen at the job, and the care, tending and human supervision required to assure that such tanks remain operational, have not been regarded as imposing unacceptable limitations. Taking into account cost, availability, transportability and time, the cryogenic fluid system has been accepted and adopted as offering the clearly superior cooling system for performing isolated freezes at unpredictable, and possibly remote, locations. "Guidelines to Good Practice in Pipe Freezing", Version PD 1, January 1990, Corrections June 1991, (C) U. of S. & D. A. W. 1990, c:pt1CvPpip, comprehensively summarized, as of the early 1990's, the state of the art of the technology of pipe freezing. The following quotes from pages 4 and 5 of the document are instructive. There are, in principle, many different ways in which the temperature of the pipe and its contents can be reduced to achieve a freeze. For small diameter pipes, freon aerosols have been used by plumbers, but, as the use of freon is now considered an environmental hazard, this practice is now obsolete. Somewhat larger diameters can be handled by the `Jet-Freezer` type of equipment that uses liquified carbon dioxide. For larger pipes, the most common technique is to attach a jacket to the outside of the pipe and to fill the resultant enclosed space with either liquid nitrogen or a mixture of solid carbon dioxide and a heat-exchange fluid. Liquid nitrogen is relatively inexpensive, readily available and it is easily transported, stored and transferred to the freezing jacket. Furthermore, its boiling point of 77 K (-196° C.) is far below the temperature required to freeze water and many other fluids and high heat extraction rates are thus possible due to this large degree of supercooling. It is also the only convenient way to obtain temperatures low enough to freeze crude and many refined oils. The circumstances in which the use of liquid nitrogen is essential are listed in the flowsheet, Factors Dictating the Use of Liquid Nitrogen. However, as noted earlier, many of the pipes to be frozen are made from ferritic steels which are well below their tough-to-brittle transition at the boiling point of liquid nitrogen. Furthermore, the direct impingement of liquid nitrogen on the pipe wall can create large temperature gradients and thermal stresses in the pipe wall. For these and other reasons, it is often necessary to use a different cooling technique, the so-called controlled temperature freeze. The flow sheet entitled Factors to be Considered for Use of a Controlled Temperature Freeze addresses these points. In a typical controlled temperature freeze an intermediate heat-exchange fluid, such as isopropanol or methanol, is used to surround, and be in contact with, the pipe wall. This fluid is in turn cooled either by the direct addition of solid carbon dioxide, by the controlled supply of liquid nitrogen, or by the circulation of a refrigerant through a coil immersed in the fluid. Cooling times are much longer using this technique, but fine control can be maintained over both the cooling rate and minimum temperature. It is thus the preferred technique for use in suitable applications where possible brittle failure of the pipe is an unacceptable risk. However, as most of these heat-exchange fluids are flammable, their use increases the fire hazard and this risk may not always be acceptable, particularly in sensitive locations. Alternative inert heat-exchange can sometimes be substituted, but they limit the usable temperature range. Recent developments in Remote Coolant Circulation, in which the working fluid is cooled using a heat-exchanger coil and then circulated within the freezing jacket, may offer a viable alternative technique for controlled temperature freezes. Once the general concept of pipe freezing have been grasped, it is necessary to decide whether is it likely to be possible to carry out a successful freeze under the prevailing circumstances. In all except the most straightforward cases, it is highly desirable for a potential contractor to be able to carry out an on-site survey to establish the relevant details. For freezes on off-shore oil platforms and other complex installations a site survey is essential. The flowsheet, Factors to be Considered to Determine Whether a Particular Freeze Can be Carried Out, addresses many of the points that have to be considered at this stage. The first, and relatively obvious, point is that there must be reasonable access to the site at which the freeze is to be effected. It must also be possible to bring adequate supplies of cryogen onto the site, to re-supply when necessary and to be able to transfer the cryogen to the point of freeze. In the flowsheet, some of these factors have been given a `Star` rating to signify their importance. If any of these factors are at their unfavorable bound it is most probably going to be difficult, if not impossible, to achieve a successful freeze. The above quote makes clear that within the art of pipe freezing there is a distinction between technology suitable for freezing "small diameter" pipes and technology applicable to freezing "larger" pipes. The present invention falls within the art of freezing "large" pipes. The convention is adopted herein, commensurate with general industry understandings, that "large" pipes comprise pipes of 6" outside diameter and up. Such pipes may well have an outside diameter ("OD") of 30" or 40" or more. The above quote also illustrates that as of the early 1990's the industry practiced some variation of the cryogenic fluid system for freezing "large" pipes. "Subsea" pipe freezing forms a relatively new arena and area of expertise. The first subsea pipe freeze experiments and research appear to have surfaced in the mid- to late 1980's. Of the only two known and reported successful commercial subsea pipe freezes, both utilized cryogenic fluids systems. One successful freeze was reported in the Offshore Engineer Scotland Report of July, 1994 and published by presentation and handout at the Society for Underwater Technology Conference at Thainstone, Scotland at the end of the summer of 1994 (herein incorporated by reference). The other successful freeze was published in "Technology" and was based on a presentation to the Energy and Environmental Expo 1995, Houston, January 29-February 1. In each case, the cryogenic fluid tanks themselves were located on the surface vessel, with insulated hoses used to pump the "cooling medium" from the vessel to a jacket located around the subsea pipe. Hoses also returned the cooling medium to the surface. "Consumed" cryogenic fluid was released to the atmosphere. The "cooling medium" comprised either the cryogen itself or a secondary coolant working fluid that was circulated through a heat exchanger on the vessel in cooling contact with the cryogen. The source of cooling power in large pipe freezes has historically been located where it is readily accessible for human handling and management. In particular, cryogenic fluid tanks require a significant amount of constant human attention and handling to keep them operational. In contrast to industry practice, the present invention teaches away from the cryogenic fluid system, and in particular, away from the extensive use of insulated hoses. The method and apparatus herein disclosed for freezing large pipes, including subsea applications, is founded, in fact, upon subsea experience. A plug must be maintained frozen during the period of pipe modification and repair. Loss of a plug after the pipe has been cut could be hazardous. Insulated hoses for the transport of the cooling medium can sustain only a limited amount of flexing without danger of fatigue or over stressing under tension, due to ice build-up from the cryogen. Therefore, rough weather could require a vessel to pull off all hoses transporting a cooling medium to a subsea pipe freeze to avoid loosing hose integrity. If pulling off were required after the pipe was cut, the potential is raised that the plug and a pipeline of fluid could be lost before insulated hoses could be reattached in calmer weather. The inventors have determined, based upon the above considerations, among others, that the better practice in pipe freezing adopts a cooling methodology wherein (1) the source of cooling power is stationed "proximate" to the section of pipe to be frozen. This method avoids the necessity for maintaining extensive insulated hoses transporting a cooling medium and the limitations and dangers such hoses impose, especially in a subsea environment. "Proximate" will be used somewhat arbitrarily herein to mean at least within about twenty feet. Ten feet or less would be preferable. Ten feet or less will be referred to herein as "adjacent". The above teaching implies that if the pipe is subsea, the source of cooling power should be stationed subsea. It can be added that, not only may insulated hoses transporting refrigerant from the surface to a sea bed present unacceptable environmental risks, they present occupational safety issues and possibly prohibitive costs as well. If the realistic ability to ensure the integrity of insulated hoses throughout the time of a subsea pipe freeze is inadequate, then the whole operation may be jeopardized. Further, such hoses are expensive and their capacity to be reused from job to job is doubtful, due to the complex loading experienced in offshore operations. Inability to reuse insulated hoses significantly raises their costs. The hoses must be regarded as "consumed" with each job. The present invention, therefore, teaches accepting, as a limitation, the avoidance of transporting refrigerant through insulated hoses over significant distances, and in particular, from surface to subsea. Accepting such a limitation has led the inventors to teach placing the source of refrigerating power adjacent to, or proximate to, the section of pipe to be frozen, including on the sea bed or on the subsea pipe, when applicable. A second aspect of divergent methodology taught by the present invention follows the acceptance of the above limitation, namely that the source of refrigeration should be placed proximate to pipe. Acceptance of that first limitation has led the inventors to rethink and reject the industry standard cryogenic fluid system. The present invention also teaches, therefore, what might appear to the industry initially as improbable, namely, (2) the use of a recycling refrigeration system for large pipe freezes, including subsea pipe freezes. Use of a recycling refrigeration unit eliminates the use of cryogen tanks and the venting of used gases. As an associated innovative aspect, the present invention further teaches (3) using the encircling pipe jacket itself for the evaporation chamber of the refrigeration unit. Such design enhances cooling efficiency and distribution without sacrificing temperature control. Such design eliminates the inefficiency and complications associated with use of secondary "working" fluids in order to perform a "controlled temperature" freeze. At least for large pipe freezes, the concept of (1) using a recycling refrigeration unit with (2) the jacket encircling the pipe configured as the unit's evaporation chamber, as well as the concept of (3) stationing such cooling unit proximate the pipe subsea, is not within the teaching of the industry. This is abundantly clear as of the early 1990's from the "Good Practices in Pipe Freezing" document. The present applicant itself performed one of the two above reported known and commercially successful pipe freezes. (The North Sea freeze--see document incorporated by reference). As of the beginning of 1994, it is clear that the applicant itself did not conceive of the need for, nor the value of, nor the possibility of, the method and apparatus comprising the present invention. Brister, in U.S. Pat. No. 4,441,328, may have appreciated the value of avoiding the use of extensive insulated hoses running subsea for large pipe freezes. Brister suggests the possibility of locating a cryogenic fluid system in a seabed habitat, or upon a submersible platform. A habitat, however, of a size and scope where humans could maintain the cryogenic tanks in working order, (clearly suggested by Brister because cryogenic tanks require a large amount of care and tending,) involves such expense and logistics as to render Brister's suggestion impractical. Over a several day operation many tanks would need to be changed out and used. Likewise, submersing cryogenic tanks without ready human supervision in a subsea environment is impractical. The present invention, in addition, includes second and third aspects, e.g. a recycling refrigeration unit with the jacket as evaporator, that were not at all taught, suggested or appreciated by Brister. It has been commonly believed that sufficient cooling power to effect a timely freeze in a large pipe required the high cooling power of a cryogenic fluid system. The disclosure of a sufficiently mobile yet powerful recycling refrigeration unit design is unanticipated. Part of the refrigeration unit design avoids the inefficiency of "secondary working fluids" for performing "controlled freezes" by incorporating the pipe jacket as the evaporative chamber, a point also unappreciated and untaught by Brister. One patent and two applications, one application perhaps never filed and still secret, have been identified that teach the use of a recycling refrigeration system in connection with pipe freezing. The one patent, Radichio, not surprisingly, dates back to the 1970's, before industry had accepted and adopted the superiority of the cryogenic fluid system. None of these documents teach or suggest the capacity of locating a recycling refrigeration unit subsea. None of the documents address how the refrigeration unit is physically connected to the jacket or the subject matter of extensive insulated hosing. None of the documents teach a refrigeration unit design wherein an encircling pipe jacket suitable for large pipe freezes forms the evaporative chamber and thereby avoids the inefficiency of using a secondary working fluid, without loss of a "controlled" freeze. A German Patent application paper or draft (filing information unknown) that has come to hand discloses a topside recycling refrigeration system that requires, and calls for, a secondary working fluid to be piped to the jacket. Among other points, the application does not appreciate the possibility of the design, nor the efficiency nor control that is possible with, implementing the pipe jacket as the evaporative chamber. Radichio, in U.S. Pat. No. 4,309,875, first of all, does not appear to address the problems inherent in the "large" pipe freeze art. One of skill in the art would understand Radichio as dealing with small pipe and low pressure problems in a topside scenario. Radichio does not conceive of, teach or suggest performing a pipe freeze utilizing a jacket for encircling a pipe as the evaporative chamber. Radichio teaches a cradle evaporative chamber. The cradle is taught to never come in contact with the pipe. Like the German application, a secondary cooling medium is always interspersed between the cooling cradle and the pipe. The plug formed from Radichio's cradle would be asymmetrical. Supplying sufficient cooling power while minimizing the size of the refrigeration unit and producing a sufficiently low pipe surface temperature to freeze a plug capable of withstanding substantial pressure are clearly not issues of concern for the Radichio design. Freezing a line by means of a cradle would not optimize or economize the cooling power expended. Radichio shows neither motivation to, nor capability to, solve the issue of optimal evaporative chamber design for large pipe freezes. Clearly, for large pipes it is unlikely that a cradle system could provide sufficient cooling power to close the plug. The Mark/Shell Patent Application GB 2 195 738 A suggests that jackets encircling the pipe might be provided with a cooling circuit which might consist of a self-contained compressor driven refrigeration system. The published Application, subsequently abandoned, contains no illustration nor reduction to practice. The Shell document is ambiguous or silent as to whether a secondary cooling agent is intended and/or whether extensive insulated hoses are necessary. It does not teach or suggest landing a refrigeration unit, or any cooling unit, proximate the pipe. SUMMARY OF THE INVENTION The invention comprises apparatus for freezing plugs in large pipes, including subsea pipes. The apparatus includes a jacket adapted to encircle a section of large diameter pipe. A recycling refrigeration unit is connected to the jacket, preferably located proximate the jacket, and incorporates the jacket as an evaporation chamber. For subsea environments a subsea submersible housing encloses a portion of the refrigeration unit, and the housing may be attached to the pipe. The large pipe is anticipated to comprise pipe of outside diameter between 10 and 50 inches. Preferably, the jacket has a length of at least one and one half times the pipe outside diameter and defines within it a baffled interior fluid passageway. In preferred subsea embodiments the condenser is attached outside of the submersible housing and the refrigeration unit includes means for reversing its cycle to heat the jacket. It is anticipated that seawater will infiltrate and occupy at least portions of space between the jacket and pipe when attached to the pipe. In some embodiments the refrigeration unit includes a single stage rotary compressor and a sub-cooling cycle. The method for freezing a large pipe plug includes surrounding a portion of large pipe with a jacket wherein the jacket comprises an evaporation chamber for a proximately located recycling refrigeration unit. Refrigerant is cycled through the unit and the jacket. Preferably, means exist for reversing the cycling of the refrigerant to the unit and the jacket. In subsea environments the method includes landing a recycling refrigeration unit under seawater proximate a portion of pipe, and preferably attaching the unit to the pipe. BRIEF DESCRIPTION OF THE DRAWINGS A better understanding of the present invention can be obtained when the following detailed description of the preferred embodiment is considered in conjunction with the following drawings, in which: FIG. 1 provides a process schematic for a preferred embodiment for a refrigeration system of the present invention. FIG. 2 illustrates a preferred embodiment for a subsea refrigeration system, the illustration providing an end elevational view. FIG. 3 provides a side elevational illustration of the embodiment of FIG. 2. FIG. 4 illustrates structural elements of an evaporator jacket. FIG. 5 illustrates structure and flow within an evaporator jacket. FIG. 6 illustrates a preferred embodiment for a surface refrigeration system of the present invention. FIG. 7 comprises a diagram of a simplified vapor compression refrigeration system. FIG. 8 comprises a diagram as in FIG. 7 but including rough estimates of ideal pressures and temperatures. FIG. 9 comprises an end elevation schematics of an evaporator jacket, illustrating hydraulic actuators. FIG. 10 illustrates schematically a refrigeration unit of a preferred embodiment with indication of location of parts with respect to a subsea pod housing. FIG. 11 illustrates possible structure for an evaporator pipe jacket. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS This invention relates to the formation of a temporary isolation in a subsea, topside or on-land pipeline system used to transport products including oil, gas and water. The method described comprises a novel version of a pipe freezing technique where a liquid such as water, oil or gel is frozen within a pipe to form a plug. Traditionally, large pipe freezing has been performed using methods where a fabricated enclosure or jacket is placed around a pipe to form a chamber through which a cryogenic liquid, generally liquid nitrogen, is pumped to cool down the surface of the pipe and ultimately form a frozen plug within the pipeline. This technique has been demonstrated both in air and subsea environments and has been further extended to allow freezing at controlled temperatures by using a heat exchanger to pre-cool a working fluid refrigerant which is then itself circulated around a jacket. All known subsea freezes have required the utilization of significant lengths of insulated hoses for transporting a cooling medium from the surface to the subsea pipe. The present invention provides a new method for the formation of ice-plugs in large pipeline systems to allow, for instance, intervention in the pipeline for the purpose of repair or modifications. It is recognized in such cases that the ice plug must form and hold, against pipeline and ambient pressure, a temporary barrier against loss of line product and, in the case of subsea pipeline, ingress of sea water. The novel method and apparatus of the present invention involves the use of a potentially remotely controlled but proximately located recycling refrigerator system where the evaporating chamber of the system is constructed in such a way as to encompass the pipeline at the designated area. Such an encompassing chamber may be referred to as a jacket. The system is suitable for use both subsea and on the surface. In the case of subsea operation, the refrigeration unit is fully submersible to the seabed. Control of the submerged refrigerator system is preferably provided from the surface, e.g. from a vessel or platform or possibly a submersible unit, via a control umbilical which can provide power, communication and monitoring between the freezing system and the surface. A flexible umbilical system can tolerate substantial vessel heave in bad weather conditions. For surface or topside application of the system, the refrigerator is man portable to facilitate equipment set up. The invention provides a system for forming ice plugs in large diameter pipeline systems, or in a pipe of six inches outside diameter and larger. The plug is formed in a liquid, which may comprise the line product, an alternative liquid injected into the pipeline such as gel, or a slug of liquid which has been displaced through the pipeline to a pre-determined position. Both formation of the plug and control of the equipment and monitoring is possible from a location remote from the freeze site. The system is based primarily on a recycling refrigerator that is suitable for both subsea and surface use. The principle of operation of the unit, whether subsea or surface, is the same. The housing and the configuration of the system may be specific to the application. FIG. 1 shows a schematic of the principal components of a refrigeration system in accordance with the present invention. High pressure, high temperature liquid refrigerant L flows through liquid line 1 to a sub cooler 2, which preferably, in some designs, pre-cools the liquid refrigerant to increase the efficiency of the system. A sub cooler is optimally employed with a single stage compressor. If a specific site application permitted, a multiple stage compressor might be preferred. It might be mentioned that in order to improve the performance of some types of vapor compression refrigeration systems, especially with a single stage or rotary compressor, it is a common practice to pre cool the refrigerant liquid prior to injection into the evaporator. This can be accomplished in several ways. In one preferred embodiment of the subsea system a direct expansion heat exchanger is used, e.g. sub cooler 2. The principle advantage of sub cooling is, in effect, to reduce the mass flow of the refrigerant through the evaporator per unit of heat extraction. This is due to the reduction of "flash gas" (refrigerant evaporating to cool remaining liquid refrigerant to the evaporator saturation temperature) at the refrigerant control device. Subsequent to subcooler 2, liquid L flows to refrigerant flow control valves 3 and 4, which may be thermostatically or electrically controlled. As the liquid passes through the control valves, the pressure is reduced to that of the evaporator 5 so that the saturation temperature of the refrigerant entering the evaporator will be below the temperature of the refrigerated item i.e. the pipe. In the preferred embodiment illustrated evaporator 14 is incorporated into jacket J. Part of the liquid vaporizes as it passes through the flow control valves in order to reduce the temperature of the remaining liquid to the evaporating temperature. In the evaporator the liquid vaporizes at constant pressure and temperature as heat (to supply the latent heat of the vaporization) passes from the pipe, through the evaporator walls and to the vaporizing liquid. By action of a motor driven rotary compressor 6, vapor V resulting from vaporization is drawn from the evaporator, through suction line 7, and into the suction inlet of the compressor. Vapor V leaving the evaporator is saturated, and its pressure and the temperature are those of the vaporizing liquid. In the compressor the temperature and pressure of the vapor are raised by compression and the high temperature, high pressure vapor is discharged from the compressor into discharge line 8. The vapor flows through the discharge line into coaxial condenser 9 where it gives up its heat to the condensing media. As the hot vapor gives up heat, its temperature is reduced to the saturation temperature corresponding to the discharge pressure, and vapor condenses back to liquid as further heat is removed. By the time the refrigerant leaves the condenser all the vapor is condensed back into the liquid phase. The liquid L is then re-circulated. A motor cooling system 10 may also be included to reduce the ambient temperature of the refrigerator system housing, or enclosure, if required. Such requirement is more likely, for example, in warm Gulf of Mexico waters than in cool North Sea waters. The refrigeration system above described preferably can be designed and run in a heat pump mode (reverse cycle), as is known by those in the art, in order to provide the capacity to heat the jacket for rapid thawing of the ice-plug and/or for removal of the system from the pipe upon completion of an operation. Alternatively, electrical heaters may be incorporated on the jacket assembly (not shown) for the same purpose. In the case of a subsea system (see FIGS. 2, 3 and 4) the refrigerator plant, i.e. the motor, compressor, subcooler, motor cooler, and valving, are enclosed in a pressure vessel or enclosure 11, alternately referred to as a housing or subsea pod, for submersion to water depths up to and in excess of 300 m without ingress of water. The enclosure may be fixed to a support frame or saddle 12 which supports the full weight of the system and attaches to and transfers the load through buffer blocks 13 directly to the pipe. Straps or other securing devices suitable for operation by divers or a Remotely Operated Vehicle (ROV) may be used to attach the refrigerator system to the pipeline. The evaporator 14 forms, or forms part of, the pipe jacket J and may be split into segments, typically three, a top section and two hinged side sections. The evaporator may be further split into a number of circuits to ensure good distribution of refrigerant and even cooling (see FIGS. 4 & 5). Non-collapsible insulation 15, such as vinyl insulation, is preferably attached to the outside of the evaporator to minimize heat gain to the freeze area from the surrounding water. The jacket assembly can incorporate instrumentation to monitor the pipewall temperature, jacket temperature and an ice plug closure detection system which may be an acoustic based system. Models predicting the growth and condition of the plug can be developed based upon acoustic sampling adjacent a freeze jacket. The refrigerator system can include hydraulic actuators to facilitate jacket installation and can be designed in such a way as to facilitate installation by both diver or ROV. The condenser 9, FIGS. 2 and 3, in a subsea system may be located on the outside of the enclosure, and thus submerged in water. Cooling can optimally be provided from the surface by pumping water down a flexible hose to such condenser and dumping the heated water subsea. Alternatively, using an alternate condenser design, heat could be dumped directly to the surrounding sea water without the requirement for a water supply from the surface. The compressor drive motor could be driven by 3-phase electrical supply from a surface vessel or platform. Back up submerged batteries, including subsea batteries, can be provided in case of the necessity of pulling off a surface vessel in bad weather. Such batteries could maintain a frozen plug during the workover operation. A multicore umbilical preferably would be run from the surface to the subsea refrigerator system to provide both power, control and instrumentation lines. Remote control of the system could include start up and shut down of the refrigerator, motor speed adjustment and valve actuation. In addition to jacket and pipewall instrumentation, temperature and pressure monitoring of the system can be provided via the same umbilical. The umbilical itself can incorporate a wet-mateable connector 16 to allow for umbilical replacement subsea in the event of failure, and would be designed to be flexible to allow for vessel heave, even in extreme weather conditions. Alternatively a hydraulic motor could be used to drive the compressor. In this case, the exhaust from the motor could provide a condensing media. On the surface, a control unit can be employed to operate and monitor the deployed refrigerator system. A data logging system may also record all instrumentation data throughout the whole operation. The control unit may also operate, optionally, as a "closed-loop" system where no operator input is required to control the cool down or the maintenance of the pipe during freeze. Software control could not only, for example, in conjunction with the refrigeration unit of the present design, maintain the pipewall at a specified temperature but could also monitor various parameters within the operating unit to allow automatic control of housing environment temperature and refrigerant sub-cooling. For surface or topside use the refrigerator would in most cases be located off the pipe (see FIG. 6). Generally, it would not be practical, due to space restriction, to mount the refrigeration system directly onto the pipe. However, the refrigeration unit of the present design is sufficiently compact and mobile as to be allowed to be placed proximate the pipe to be frozen in even those locations where space is scarce, such as in plants and offshore platforms and on vessels. The top side unit evaporator 17 would again be designed to encompass the pipeline. For smaller diameter "large" pipe such evaporation design might comprise a single flexible piece which could accommodate a range of pipe sizes. For larger "large" pipe sizes the evaporator would preferably comprise multiple segments, as illustrated in the subsea system. With the subsea system, a thin layer of ice is expected to form between at least portions of the evaporator jacket and the pipe. Although the jacket is preferably designed to closely encompass a pipe, jacket design will allow for variation in actual pipe diameter in practice. Thus, portions of the jacket may directly contact the pipe and other portions may not, but means for securing good thermal contact can be provided to accommodate these pipe size variations. Ice filling the space between the jacket and the pipe forms one means of supplying a good conductive heat path. In the case of a surface system, where the pipe is surrounded by air a heat transfer contact compound might be provided. As a further topside alternative, a chamber could be attached around the evaporator itself to allow submersion of the jacket in a liquid such as water. After installation of the evaporator, preferably an insulation 18 is fitted adjacent the jacket to prevent heat transfer from the environment to the freeze area. Whereas with a subsea system it is desirable for suction and discharge lines to a jacket to be hard-piped, with a surface unit it is more practical for suction and discharge lines 19 to be flexible. For larger pipe diameters, separate suction and discharge lines would be used. For smaller pipe diameters, where the volume of refrigerant is comparatively small, a single coaxial flowline could be used for both suction and discharge. The suction and discharge lines, or coaxial suction and discharge line, of the refrigeration unit may be fitted with connectors to allow for disconnection of the evaporator in order to design for large variations in pipe diameter freezes with a single refrigeration unit. Such key elements and system make-up for a surface unit could be essentially the same as those described for a subsea system. However, an air/gas cooled condensing unit 20 may be preferable topside, as an alternative to a liquid cooled unit, and either a single stage rotary or multi-stage compressor might be used topside. The drive motor for a compressor could be either electrically, pneumatically or hydraulically driven. The topside surface unit housing (or enclosure) 21 could contain all of the refrigerator system components with the exception of the evaporator. The topside housing is designed such that it is light and can be man-handled in areas of restricted access, such as offshore platforms and industrial plants and vessels. The instrumentation on the topside system would provide similar operating and monitoring capabilities as with the subsea system. The topside housing could accommodate data readouts, including motor speed, section and discharge temperatures, pipewall temperature and other critical parameters to allow a determination of the refrigerator system performance. An interface to allow external data logging/recording could be incorporated. The control system on the refrigerator unit may also operate as a closed-loop system where no operator input is required to control the cool-down and maintenance of the pipe during the freeze. The possible advantages of the system of the present invention, for both a subsea and surface environment, can be summarized as follows: Subsea Unit 1. Elimination of significant insulated hoses exposed to environment and vessel heave. 2. Elimination of lengthy equipment set-up time. 3. Elimination of need to perform lengthy drying/purging of insulated hoses and jacket following subsea installation to pipe. 4. Vast reduction in cost of consumables, such as cryogenic fluid and insulated hoses. 5. Significant reduction in quantity of equipment and space requirements. 6. Minimization of need for equipment operators. 7. More reliable: less components (less moving parts). less chance of operator error. 8. Little operator input required. Prior art systems require continual operator control. 9. Better control over pipewall temperature. 10. Option to speed up thaw time by using reverse cycle or heaters. 11. Reduction in overall time to perform freeze. Surface System 1. Reduction in equipment set-up time. 2. Equipment man portable. 3. Reduction in hoses running in work areas. 4. Size of equipment allows it to be sited proximate the freeze location. 5. No consumables, so reduction in cost. 6. Significant reduction in area required by equipment 7. No need to purge/dry equipment. 8. Reduction in number of operators. 9. Little operator input required. 10. Better control over pipewall temperature. 11. Option to speed up thaw time by using reverse cycle or heaters. 12. Reduction in overall time to perform freeze. A diagram of a simplified vapor compression refrigeration system is shown in FIG. 7, wherein the evaporation unit comprises a jacket whose temperature can be controlled. The principle parts of the system comprise an evaporator pipe jacket J, whose function it is to provide heat transfer service through which heat can pass from the refrigerated item into the vaporizing refrigerant. A suction line S conveys the low pressure vapor V from the evaporator J to the compressor CP inlet. A vapor compressor, whose function it is to remove the vapor from the evaporator and to raise the pressure and temperature of the vapor to the point where it can be condensed with normally available condensing media is connected to the suction line. A discharge line D delivers high pressure, high temperature vapor V from the discharge of the compressor CP to the condenser CD. Condenser CD provides a heat transfer surface through which heat passes from the hot refrigerant vapor to the condensing media. A receiver tank R provides storage for condensed liquid refrigerant L. A liquid line LL carries liquid refrigerant from the receiver R to the refrigerant flow control FC. A refrigerant flow control meters the proper amount of refrigerant to the compressor to reduce the pressure of the liquid entering the evaporator so that liquid will vaporize in the evaporator at the desired low temperature. The unit will be designed and scaled, including the selection of the refrigerant, to operate within a jacket temperature range. The proper jacket temperature range will be selected in accordance with the lower limit to which the pipe temperature can be safely reduced. Such a refrigeration system can be conceptually divided into two parts according to the pressure exerted by the refrigerant in the two parts. The low pressure part of the system consists of the refrigerant flow control, the evaporator, and the suction line. The pressure exerted by the refrigerant in these parts is the low pressure under which the refrigerant is vaporizing in the evaporation. This pressure is known variously as the "low side pressure", the "evaporator pressure", the "suction pressure", or the "back pressure". The high pressure side of the system consists of the compressor, the discharge line, the condenser, the receiver, and the liquid line. The pressure exerted by the refrigerant in this part of the system is the high pressure under which the refrigerant is condensing in the condenser. This pressure is called the "condensing pressure", the "discharge pressure" or the "head pressure". The dividing points between the high and low pressure sides of the system are the refrigerant flow control, where the pressure of the refrigerant is reduced from the condensing pressure to the vaporizing pressure, and the discharge valves in the compressor, through which high pressure vapor is exhausted after compression. Roughly estimated pressures and temperatures relating to a basic vapor compression refrigeration system are shown on FIG. 8. Such roughly estimated pressures and temperatures are "ideal". Pressure drops within the lines or heat exchanges are not considered. In actual systems such pressure drops can be substantial. Additional control and subsystems can enhance the performance and operating range of the simplified system illustrated above, and may be incorporated into preferred embodiments of the present invention. Such additional controls and subsystems include the design of an evaporator which forms part of the pipe jacket. For larger pipe, this evaporator is preferably split into three pieces, a top section and two hinged side sections. See FIG. 9. The evaporator may be further split, in a preferred embodiment, into four circuits to insure more even distribution of the refrigerant. Two suction lines can be used to convey gas from each evaporator side section. A single stage rotary vane compressor may optimally be selected for some preferred embodiments, taking into account refrigeration unit size limitations and the limitations upon the temperature to which the pipe can be safely reduced. Such a compressor has an additional inlet commonly known as an economizer port 25. See FIG. 10. The economizer port comprises effectively an intermediate pressure port. This intermediate pressure port is a port where additional gas can be introduced into the compressor part way through the compression cycle. Such offers the advantage that, in some respects, the compressor behaves like a two stage machine. The feature can be used to provide additional subcooling to the refrigerant liquid, enhancing system performance. Preferred embodiments may further include a discharge line, a water cooled coaxial condenser and a receiver which in some field versions may be omitted, as the condenser tube size volume could be selected so as to act both as condenser and receiver. A liquid line and four refrigerant flow controls are preferably included, two controls for the jacket, one control for the subcoolant and one for the motor cooling system. Automatic thermostatically controlled valves may be utilized. The jacket valves may be replaced with electronic valves in some embodiments. A subcooler may comprise a plate type heat exchanger that can be used to reduce the temperature of the liquid refrigerant prior to introduction into the evaporator. Use of a subcooler has the effect of refining mass flow through the evaporator for given evaporator duty. An optional motor cooling system comprises a small evaporator, expansion valve and temperature regulator. The compressor drive motor may include a small impeller attached to the drive shaft. The drive shaft can draw fluid through the evaporator and over the motor body. The regulator throttles a suction line from the evaporator when the temperature falls to a preset level. An oil cooler may be included in order to keep discharge temperature to acceptable levels. In preferred embodiments the refrigeration unit may be operated in heat pump modes, which means reverse cycle, as is understood in the industry, in order to heat the jacket for rapid thawing of the ice plug. Alternately, the jacket may incorporate electrical heaters into the assembly. The compressor of preferred embodiments may be driven by a hydraulic motor. When hydraulically driven, the exhaust from the motor may be used as condensing media. The evaporator of preferred embodiments may be constructed from flexible lines in order to fit irregular shapes. The evaporator assembly may also incorporate pipe wall temperature sensors SS, connected by line SL to a unit control system and/or acoustic sensing devices AS. See FIG. 2. Acoustic devices sense and transmit the sounds made by forming ice crystals. The pattern of sounds can be correlated with plug formation model data to sense the growth and condition of the plug and predict the stages of the formation of the plug. The jacket may be so constructed as to facilitate installation by remotely operated vehicles. Automatic release of subsea umbilicals may be incorporated into the design. The refrigeration plant may be positioned off the pipe and connected by flexible lines, while retaining the proximity of the refrigeration unit to the jacket. The jacket may be designed so that it can be closed by hydraulic actuators 26 if necessary. See FIG. 9A. The method and apparatus of the present invention is designed to freeze pipes of up to at least approximately 48 inches in diameter, and larger, achieving a high rate of heat transfer. A rate of heat transfer is primarily a function of pipe diameter, ice thickness and jacket delta temperature. The method and apparatus is capable of operating at depths in excess of 1,000 feet. It affords simple and practical operation and installation with an ability to accurately control pipe surface temperature, such as to a minimum temperature of -50° C., or to -65° C., as examples. Such achieves an economy by eliminating the complication and inefficiency of a secondary working fluid. A further advantage of the invention comprises low capital and operating costs. As illustrated in FIG. 10, the method and apparatus of an embodiment may include motor M. Motor M is preferably hydraulic. A seawater supply can be used to power the motor, and the motor can vent to the sea. The motor will be located within the subsea pod H comprising the submersible housing. The apparatus may further include a hydraulically driven condensing unit CD within the submersible housing. The condensing unit connects to an exterior pipe jacket J functioning as the evaporating chamber of the recycling refrigeration unit. In preferred embodiments the condenser of the recycling refrigeration unit is located exterior to the submersible housing, thereby being cooled by the ambient water. A controlled system will communicate with remotely located users to control and maintain the temperature of the pipe surface. The control system and preferred embodiments may include an electrical umbilical connector. In one embodiment the refrigeration unit comprises a water cooled hydraulically driven condensing unit running on a refrigerant such as Forane FX 10 or R 502. In order to minimize the size of some units, a single-stage vane type compressor CP with economizer port 25 may be selected. A remote hydraulic pump, such as on a surface vessel, may supply the condensing unit drive motor with high pressure water, such as sea water, via a single hose HO. See FIG. 10. No fluid return line is necessary as the sea water can be vented directly at the pod. The hose can be sufficiently flexible to withstand significant vessel heave and be reusable. Pipe jacket temperature can be accurately controlled by either varying the pump flow rate, which determines compressor speed, or by limiting the evaporating pressure with the regulator. The apparatus of the present invention offers an environmentally friendly system with no pollution potential nor unusual hazard to personnel. The low thermal mass and high efficiency of the system result in a rapid jacket pull-down, saving time. The system is easier to install on pipe than other industry common systems since the system does not entail bulky hoses. The supply hose HO of the present system would comprise approximately only 10% of the cost of the composite insulated hoses currently employed in the cryogenic industry. Furthermore, only half as much hose would be required, as no return line would be necessary and the hose is reusable. The present invention makes it practically possible to provide a 100% backup at the location for all active refrigeration components. Power backup can be provided with batteries. Furthermore, the apparatus of the present invention can be designed to operate at depths in excess of 1,000 m. Preferred embodiments of the present invention may also utilize thermal panels for the construction of the jacket. Thermal panels provide a highly developed heat transfer surface commercially available in a variety of double and multi-wall forms. The high surface-to-volume ratio of thermal panels result in preferred heat transfer characteristics. FIG. 11 illustrates possible and commercially available configurations, design and structure for evaporator pipe jacket J in order to maximize heat transfer and heat transfer surfaces. FIGS. 4 and 5 illustrate other possible embodiments for the interior structure of a jacket J, embodiments defining an interior baffled path for the liquid/vapor. It is preferable that the liquid/vapor L/V pursue a circuitous route through the jacket comprising the evacuation chamber to optimize heat transfer. FIG. 5 illustrates baffling B in evaporator 14, baffling B being indicated by dashed lines, which baffling defines a circuitous route in FIG. 5 for the refrigerant. The foregoing disclosure and description of the invention are illustrative and explanatory thereof. Various changes in the size, shape and materials as well as the details of the illustrated construction may be made without departing from the spirit of the invention.

Apparatus and method for freezing large pipe plugs, including in particular subsea applications, wherein the apparatus includes a jacket adapted to encircle a section of large pipe, a recycling refrigeration unit having said jacket connected as an evaporation chamber and, for subsea applications, a submersible housing enclosing a portion of such refrigeration unit; the method including landing a submersible recycling refrigeration unit on or proximate a pipe and cycling, and preferably reverse cycling, a refrigerant therethrough.

BACKGROUND OF THE INVENTION This invention relates to wastewater treatment and in particular to a process for the anaerobic biological purification of wastewater containing organic substances as well as to an apparatus for conducting the process. As is known, organic wastewater ingredients are metabolically degraded during anaerobic wastewater purification in a succession of reaction steps to methane and carbon dioxide. (For more deails, reference is directed to "Wastewater Engineering: Treatement Disposal Reuse", Metcalf & Eddy Inc., Revised by Tchobanoglous, 2nd Edition, 1979, Boston). In this connection, the conversion of the most slowly metabolizable substances determines the total reaction period. In case of wastewaters having a complex composition, wherein part of the organic load, expressed as COD, exists in a difficult to metabolize form, such as, for example, as undissolved solids, polymers, or polycondensed aromatics, this leads frequently to long reaction times and large-volume reactors. A great variety of different anaerobic reactors are presently in use, such as, for example, complete-mix, single pass reactors, anaerobic reaction tanks with sludge recycling via a post clarification tank, upward-flow reactors with internal sludge retention, or anaerobic solid-bed reactors. Besides single-stage reactors, two-stage reactors are likewise employed with a separate acidifying and methanizing stage. For an illustration and discussion of such reactors, reference is directed to Anaerobic Digestion, Applied Science Publishers LTD, London, 1980; Anaerobic Digestion, 1981 Elsevier Biomedical Press, Amsterdam-New York-Oxford 1982. SUMMARY An object of this invention is to provide a process of the above-discussed type, as well as an apparatus for conducting the process, in such a way that a high purifying efficiency can be attained in a simple and economical fashion, with a saving in reactor volume or with a shortening of the residence time of the wastewater in the anaerobic reactor. Upon further study of the specification and appended claims, further objects and advantages of this invention will become apparent to those skilled in the art. These objects are attained according to this invention by separating the substances from the wastewater which exhibit a low rate of metabolism under anaerobic conditions into a secondary stream prior to introducing the mainstream of the wastewater to be purified into a reactor operated under anaerobic conditions. These substances in the secondary stream are treated independently of the substances which have remained in the mainstream of the wastewater to be purified and are rapidly degradable by anaerobic microorganisms. By low-rate substances, examples thereof include, but are not limited to, partially dissolved and partially macromolecular materials, e.g., proteins, long chain fatty acids, fats, vegetable oils, tallow, bacterial and yeast cell-walls, celluloses, hemicelluloses, starch, in emulsified, suspended or colloidal state as discharged e.g. from slaughterhouses, dairies, rendering plants, oil mills, pharmaceutical and organochemical plants, pulp- and paper mills. In general, the low-rate substances are typified by a rate of metabolism which is significantly lower than rapid rate substances. For example, acetic acid, as contained in condensates of sulfite pulping plants or glucose, as contained in sugar factory wastewaters, are rapid-rate substances, for as low molecular, polar substances they are readily dissolved in water, and the metabolic pathway to methane is short. Proteins, fats, vegetable oils etc, as contained in food producing plants, are on the other hand low rate substances, as they are of high molecular weight and/or are relatively non-polar; they are in a suspended, emulsified or colloidal state; and their metabolic pathway to methane is longer, requiring a hydrolysis and depolymerization step first. For a given organic wastewater load, e.g. expressed as COD, low rate substrates show an overall digestion rate which is typically only 5 to 30% of those found with rapid rate substrates. By separation and segregated treatment of the low-rate metabolizable substances in a secondary stream, the objective is attained that the reaction velocity of the most slowly occurring reactions, namely the conversion of the undissolved substances by means of enzymes into dissolved substances, as well as the hydrolysis of macromolecular substances, such as polysaccharides, proteins, and fats, is no longer a determining factor for the hydraulic residence time of the wastewater in the anaerobic reactor of the mainstream. Consequently, the design volume of this anaerobic reactor can be made less than heretofore, or, alternatively, an improved COD degradation degree can be obtained, with the reactor volume remaining the same. As for the treatment of the substances separated in the secondary stream, these are advantageously concentrated and subjected at least in part to substantially similar if not equal anaerobic conditions. Thus the low-rate substances can be converted, with the use of a small reactor volume, into dissolved compounds having a low molecular weight, such as sugars, amino acids, short chain fatty acids, or glycerol. Due to the increased concentration of the substances separated in the secondary stream, a smaller over-all reactor volume than used heretofore is sufficient for the treatment of the influent wastewater. Advantageously, the substances anaerobically dissolved in the secondary stream are thereafter introduced into the reactor located in the mainstream and operated under anaerobic conditions. This affords the advantage of eliminating any special treatment for the further processing of these substances, and the latter contribute toward methane generation in the anaerobic reactor of the mainstream. There is also the possibility of treating the substances separated in the secondary stream at least partially under aerobic conditions. This is expedient if only partial degradation is conducted in the mainstream anaerobic reactor with additional degradation of the main wastewater substances being conducted downstream of the anaerobic reactor in an aerobic treatment stage. The low-rate substances separated into the secondary stream upstream of the anaerobic reactor can then be fed directly to the aerobic treatment stage. It is likewise possible, as an alternative, to feed these low-rate substances into this stage only after they have been concentrated and converted, under anaerobic conditions, into dissolved substances exhibiting a low molecular weight. The separation of the low-rate substances from the mainstream is effected advantageously with the addition of precipitants and/or flocculants and/or adsorbents by mechanical methods. Suitable additives that can be used include, but are not limited to, iron hydroxide or aluminium hydroxide as precipitants, and activated carbon, bentonite, or bleaching clay as adsorbents. For mechanical separation, for example, sedimentation, flotation, filtration, or centrifugation can be provided. Since especially macromolecular substances are particularly well suited for adsorption or precipitation, this procedure has the result that indeed essentially only the readily metabolizable substances enter into the anaerobic reactor of the mainstream. Whether precipitation, adsorption or a combination of both is chosen depends on the nature of the low-rate substances and availability of the additives. For the removal of suspended substances, precipitation will be the preferred method, using iron salts or alum as precipitants, whereas emulsified and colloidal substances will be preferably removed by adsorption, by means of activated carbon, bentonite etc. An apparatus for conducting the process comprises at least one anaerobic reactor operated under anaerobic conditions, provided with an inlet for wastewater to be treated, an outlet for purified wastewater, as well as a gas discharge conduit for sewer gas (methane). According to the invention, such an apparatus is characterized in that the inlet for wastewater to be treated is associated with a feed means for precipitant and/or flocculant and/or adsorbent, as well as with at least one separating means for undissolved and/or flocculated and/or adsorbed substances; and that a treatment reactor for the treatment of the separated substances is connected to the separating means. This treatment reactor can herein likewise be designed as an anaerobic reactor, but its volume can be dimensioned to be relatively small since the separated substances can be well concentrated for converting them into rapidly metabolizable substances. On the other hand, the treatment reactor can, however, also be an aerobic reactor, for example if the mainstream anaerobic reactor is followed in any case by an aerobic reactor for the further degradation of the wastewater components. If the treatment reactor for the separated substances is a separate reactor, rather than being an already present reactor such as the aerobic reactor for the further treatment of the wastewater, then it is advantageous to connect the treatment reactor via a branch conduit to the main conduit leading into the anaerobic reactor downstream of the separating device and/or into the anaerobic reactor proper. This makes it possible to degrade the reacted substances, now present in readily metabolizable form, into methane and CO 2 in the mainstream anaerobic reactor, together with the substances introduced initially in readily metabolizable form. It is furthermore advantageous to provide a carrier material for microorganisms in the treatment reactor, since in such a case a high biomass concentration can be maintained and rapid conversion can be achieved of the low-rate metabolizable substances into low molecular weight, dissolved substances. Carrier materials in this connection are, preferably macroporous materials having open macropores of 0.1-5 mm, such as, for example, foam materials, e.g., polyurethane foam, ceramics, activated carbon, or swollen clay, since such materials provide a large surface area available for settling of bacteria, the latter being able to distribute themselves uniformly and firmly fixed thereon, and being forced into decentralized growth. The carrier material can be composed of one or several blocks of such a macroporous material, provided in the reactor as fixed installations, or of individual matter particles having a diameter of 0.5-50 mm. In the context of treating raw wastewater, the present invention is used to achieve a rapid rate anaerobic digestion of the whole wastewater stream by separating off the low-rate substances in a separately treated side-stream. Simple sedimentation or filtration of suspended matter e.g., would only be a partial solution of the waste disposal problem, as the precipitate or filtercake would have to be disposed off after a drying and stabilization step. The invention thus enables a joint treatment of rapid- and low-rate substances, without letting the low-rate wastes determine the volume of the anaerobic reactor. BRIEF DESCRIPTION OF DRAWING The attached figure is a schematic illustration of a preferred embodiment of an apparatus for conducting the process. DETAILED DESCRIPTION Wastewater to be treated is fed via an inlet 1 into an anerobic reactor 2, from which treated wastewater is discharged by way of an outlet 3, and sewer gas is removed via a gas discharge conduit 4. This anaerobic reactor can be designed as a complete-mix, single pass reactor, as an anaerobic reaction tank with sludge recycling by way of a post clarification tank, as an upward flow reactor with internal sludge retention, or as an anaerobic fixed-bed or fluidized-bed reactor. A feed means 5 for precipitant, flocculant and/or adsorbent and subsequently a separating means 6 for undissolved, flocculated and/or adsorbed substances are arranged in the inlet 1. The feed means 5 is suitably fashioned so that wastewater and additive are adequately blended together. The separating means 6 can be, for example, a sedimentation or flotation tank, a filter, or a centrifuge. The sludge removed from the separating means, containing the precipitated and/or adsorbed low-rate metabolizable substances, is conducted into a treatment reactor 7. Conversely, the liquid remaining in the separating means 6 is transferred, together with the readily metabolizable substances, into the mainstream anaerobic reactor 2 which can be designed smaller than heretofore by virtue of the separation of low-rate substances. The treatment reactor 7, which can also be designed to be relatively small in size, due to the fact that the separated substances can readily be treated in a concentrated form despite the slow rate of conversion of the solids into dissolved substances and the likewise slow rate of hydrolysis of macromolecules, is suitably operated in the same way as an anaerobic reactor. In this case, the effluent from this treatment reactor can be introduced via a bypass conduit 8 into the inlet conduit 1 to the anaerobic reactor 2 downstream of the separating means 6, or directly into the anaerobic reactor 2, without affecting the anaerobic conditions in the latter. Without further elaboration, it is believed that one skilled in the art can, using the preceding description, utilize the present invention to its fullest extent. The following preferred specific embodiments are, therefore, to be construed as merely illustrative, and not limitative of the remainder of the disclosure in any way whatsoever. In the following examples, all temperatures are set forth uncorrected in degrees Celsius; unless otherwise indicated, all parts and percentages are by weight. The numerical example set forth below is to clarify the extent of the savings in reactor volume attributable to this invention as compared with the conventional operation: The numerical example is based on purification of a highly loaded wastewater produced in a quality of 100 m 3 /day with a COD content of 20,000 mg/l (=2,000 kg/day). The organic load of the wastewater, expressed as COD, is to be composed of 50% readily degradable substances (=1000 kg COD x /day), 40% difficult to degrade (low-rate) substances (=800 kg COD y /day), and 10% nondegradable substances (=200 kg COD z /day). Degradation takes place in the individual reactors in the stationary condition of operation, which means that the microorganism populations of the various degradation stages are present in a high and constant concentration. Assuming, first order reaction kinetics for the COD degradation, the following equation results: dCOD/dt=k·COD; integrated: COD.sub.t =COD.sub.o /(1+k·t) wherein COD t =COD in the reactor effluent after a reaction period t (kg/d), COD o =COD in the reactor influent (kg/d), k=velocity constant of the first order (1/d), t=reaction period (d). For the readily degradable COD proportion COD x , the following applies, for example: k x =4(1/d), while for example k y =0.3 (1/d) is to be applied for the conversion of the low-rate COD proportion COD y into readily degradable COD x . The following residence time results from the aforementioned equation of the first order with the use of a conventional reactor with a throughput of the entire wastewater and assuming that 90% of the total COD is readily degradable: t.sub.x =(COD.sub.xo -COD.sub.xt)/COD.sub.xt ·k.sub.x =2.25d (wherein COD.sub.xt =COD.sub.z). Consequently, the required reactor volume for this case amounts to 225 m 3 . The following residence time results from the above equation for conversion of low-rate COD y into readily degradable COD x : t y =[(COD y +COD z )-COD yt ]/COD yt ·k y =13.3d (wherein COD yt =COD z ) corresponding to a required volume of 1,330 m 3 for the low-rate conversion of COD y to COD x . Using the procedure according to this invention, the low-rate COD and the COD impossible to degrade represents a 1% concentration. By precipitation-adsorption, a partial stream of, for example, 7% can be separated therefrom, corresponding to a volume of 14.3 m 3 . With a residence time of 13.3 days for converting COD y into COD x , a reactor volume is obtained of 190 m 3 . The total reactor volume is consequently composed in this case of 190 m 3 for conversion of concentrated COD y into COD x , as well as 225 m 3 for the total conversion of COD x into methane and CO 2 , resulting in a total volume of 415 m 3 and, as compared with a conventional reactor, in a saving of reactor volume of 69%. To faciliate comprehension of the above example, further information is provided as follows: The wastewater comes from a food-canning factory, and contains easily degradable organics, e.g. acetic acid from sauerkraut production, and slowly degradable ingredients, such as vegetable debris. The slowly degradable substances are removed from the mainstream by addition of a polymer flocculant aid and ferric chloride, with sedimentation of the formed precipitate in a sedimentation tank of conventional design. The precipitate is the above-mentioned 7% side-stream. As far as the anaerobic microorganisms are concerned, there are no specific requirements to be fulfilled. Anaerobic bacteria, as they are ubiquitous in municipal sludge digestion, e.g. are applied for start-up; those microorganisms that are best acclimated to the given substrate and reaction condition will prevail soon after start-up. In the mainstream reactor with methane production as the final metabolism step, a wide variety of microorganisms will be present in the neutral pH range (6.8-7.8), while in the side-stream reactor slightly acidic conditions (pH 5-6.5) will provide optimal conditions for hydrolyzing and acidifying bacteria, whereas methanogenic bacteria are not present. Reactor temperatures are in the mesophilic (20°-35°C. ) or thermophilic range (50°-65°C. ) with temperature control by conventional cooling, heating or heat exchange. The preceding examples can be repeated with similar success by substituting the generically or specifically described reactants and/or operating conditions of this invention for those used in the preceding examples. From the foregoing description, one skilled in the art can easily ascertain the essential characteristics of this invention, and without departing from the spirit and scope thereof, can make various changes and modifications of the invention to adapt it to various usages and conditions.

In the anaerobic biological purification of wastewater containing organic substances, some of which have a low rate of metabolism in anaerobic microorganisms, e.g., undissolved and/or partly macromolecular substances, the low-rate substances are separated from the wastewater into a secondary stream, e.g., by mechanical, adsorptive or precipitating means, before introducing the mainstream of the wastewater to be purified into the reactor operated under anaerobic conditions. The low-rate substances separated in concentrated form are treated, e.g., in a separate anaerobic reactor, or in an aerobic reactor situated downstream of the anaerobic reactor.

CROSS-REFERENCE TO RELATED APPLICATIONS [0001] This application claims the benefit under 35 U.S.C. 119(e) of U.S. provisional application Ser. No. 61/281,402 filed Nov. 16, 2009, which is incorporated by reference herein in its entirety STATEMENT REGARDING FEDERALLY SPONSORED RESEARCH [0002] Not Applicable BACKGROUND OF THE INVENTION [0003] Nebulizers for producing micro-droplets (i.e. aerosol) from liquid medicaments and presenting those aerosols for patient respiratory therapy are a well-known and practiced technology. A typical nebulizer design includes the basic elements of a gas inlet port, a respirable gas outlet, a liquid reservoir and a means for forming micro-droplets of the liquid within the reservoir. Early designs, such as represented by U.S. Pat. No. 3,097,645 and U.S. Pat. No. 3,762,409, both to Lester and incorporated by reference in their entireties herein, depict the basic elements of a typical constant-flow type nebulizer. Constant-flow type nebulizers create micro-droplets of liquid medicament based on an uninterrupted supply of pressurized gas coming through the gas inlet port and entraining liquid from the reservoir continually, forming a fraction of the liquid into an aerosol until such time the pressurized gas is stopped or the reservoir of liquid becomes empty. While representative simple nebulizers such as taught by Lester are capable of producing an aerosol, the efficiency of simply jetting an entrained liquid stream into a free space was not found to be adequate for creating micro-droplets of a consistent size and rate. U.S. Pat. No. 4,588,129 to Shanks, incorporated herein by reference in its entirety, addresses this consistent size and rate issue of the earlier Lester designs by further incorporating a fixed baffle having a convex target surface. In the Shanks nebulizer, a liquid entrained jet stream strikes upon the convex target surface of the baffle and the impact thereof allows for the momentum imbued within the liquid stream to mechanically act upon the stream and cause the creation of smaller, more readily inhaled micro-droplets at a higher rate. [0004] The Lester and Shanks nebulizers greatly advanced the art of aerosol formation, however, due to their continuous aerosol formation mode of operation, much of the liquid medicament formed into an aerosol was lost from the device during patient exhalation and idle operation of the device. Loss of aerosolized medicament to the environment is deleterious as there is a decrease in therapeutic value to the patient resulting from reduced dosing, as well as, contamination of the immediate atmospheric environment and inadvertent dosing of individuals not requiring treatment. Dosing variability with continuous aerosol formation nebulizers is also very high and largely affected by the physiological respiratory of the patient, thus two different patients with two different inhalation and exhalation time ratios using the same continuous aerosol formation nebulizer will receive significantly different doses. Improvements were then made to alter nebulizer performance such that the creation of micro-droplets through aerosolization occurred only when the patient being treated was inhaling through the nebulizer. Published U.S. Patent Application 2003/0136399 to Foley, et al., teaches a means for a nebulizer, which creates a constant micro-droplet aerosol within a closed chamber, which is released through operation of a valve. Published U.S. Patent Application 2002/0157663 to Blacker, et al., seeks to control aerosol production through patient inhalation completing the path from the liquid reservoir to the entrainment orifice and thereby allow liquid to entrain into the pressurized gas. U.S. Pat. No. 7,080,643 to Grychowski, et al., utilizes a gas diverter, which moves into and out of position wherein pressurized gas is directed across liquid transfer conduits and the vacuum created thereby causes liquid to be drawn through the transfer conduits and entrained into the gas flow. The aforementioned U.S. patent numbers are incorporated herein in their respective entireties. [0005] Although many of the problems of continuous nebulizers has been mitigated by various clinical practices, the nature of medications needed to be aerosolized for inhalation by patient has begun to change such that there is a greater need for control of dosing and environmental exposure. Previously, aerosolized medications were primarily aqueous solutions containing low mass concentrations of salts or other easily soluble compounds with wide allowable dosing profiles and low toxicities. A number of new medication have begun to be introduced, including some consisting of proteins and other biological material, that have much tighter allowable dosing profiles, greater toxicity risks, and greater concern of secondary exposure of un-intended individuals present during treatment due to exposure to exhaled aerosolized medication or aerosolized medication produced at some other time than inhalation. Some of these newer medications tend to have a much higher mass concentrations resulting in thicker solutions and higher viscosities. The result is that much more residual material may be caused to accumulate in or around the nozzle, which can impede or prevent the proper performance of the nebulizer over the course of treatment. Accumulation of material around the nebulizer nozzle, thus impeding performance, is a particular problem with nebulizers that are breath-actuated through means that include intermittently turning on and off gas or liquid flow in synchronization with patient respiration, due to the enhanced drying effect realized by these strategies. Dosing of these medications is many times a much more sensitive issue than older medications, thus a nebulizer that delivers medication only upon inhalation has a distinct advantage over those that run continuously, because inhalation and exhalation time ratios can vary tremendously from patient to patient, thus a breath actuated nebulizer can deliver a more consistent dose regardless of respiration pattern. Furthermore, breath-actuated nebulizers may help mitigate secondary exposure issues by insuring that aerosol is produced only during inhalation, although it is well known that patients will exhale some of the aerosolized medication that has been inhaled and a breath-actuated nebulizer by itself does not completely solve the problem of secondary exposure. Unfortunately, if nebulizer performance is degraded due to accumulation of medication in or around the nozzle, the benefits of breath-actuation can be largely offset by the degraded performance of the nebulizer. Therefore a need exists for a breath-actuated nebulizer that is less sensitive to material accumulation of large molecule medications, that is designed primarily for delivery during patient inhalation, and which has a design that lends itself to mitigation and control of secondary exposure. [0006] Many existing breath-actuated nebulizer involve electrical or sophisticated mechanical components necessary to detect patient inhalation. These devices suffer from high purchase price associated with added sophistication and the inconvenience of a re-usable component that needs to be stored, retained, set up, and cleaned with each use. Therefore a need also exists for a breath-actuated device that is simple in design and capable of being entirely made out of inexpensive parts and therefore potentially for single use and disposable. [0007] The present invention satisfies all of these referenced needs. BRIEF SUMMARY OF THE INVENTION [0008] The present invention is directed generally to a nebulizer for the formation of micro-droplets (i.e. aerosol) from liquid medicaments for respiratory patient treatment, and more specifically, to a baffled nebulizer wherein a static baffle used to form an atomized medicament is proximal to a shield which responds to patient respiration force to oscillate from an occluding position to an open flow position. An entrainment orifice within the nebulizer utilizes pressurized gas to draw in liquid medicament from a reservoir and to entrain that liquid medicament into a continuous high velocity stream. The high velocity liquid entrained jet is oriented such that an optimization point is achieved continuously at a defined distance in front of the entrainment orifice. A target surface is positioned at the optimization point of the high velocity jet such that the medicament within the stream is atomized into micro-droplets. When the shield is in a first position, the atomized medicament impinges upon the shield, losing jet momentum, which causes the micro-droplets to coalesce into macro-droplets that are non-respirable and return to the liquid reservoir for re-entrainment. Since the nozzle is continuously entraining and processing medication during this time, although not forming aerosol due the position of the shield, medication is continuously cycled through the nozzle region and there is little drying and little or no residual build up of thick and viscous medications. Thus the invention is less sensitive to mass accumulation and is more able to consistently deliver large molecule, high concentration, and/or more viscous medications. When the shield is in a second position, the micro-droplets are un-occluded and are released into aerosol outlet which then may be inhaled by the patient. The impingement shield moves between the first and second positions based upon the respiration of the patient, thus moving the impingement shield into the second position only when patient inhalation occurs, thereby preventing excessive waste of liquid medicament, improving patient therapy, mitigating the accumulation of medication around the aerosol producing region, mitigating secondary disposable, and doing so in a simple enough manner to lend itself well to a inexpensive and disposable device. [0009] A nebulizer assembly made in accordance with instant disclosure is capable of an expulsion rate of equal to or greater than 1.0 ml per minute at a gas flow rate of greater than 8 liters per minute and pressures of between 15 and 50 psig. The high performance of the impingement shield nebulizer at low pressures is significant in that conventional nebulizer compressors as used in home administrated therapy exhibit a pressure output of between 15 to 20 psig, a pressure range in which other nebulizer technologies exhibit diminished expulsion rates, thus requiring additional time of dosing and less than optimum aerosol particle size. [0010] In a further embodiment of the present invention, the nebulizer assembly is particularly adapted to control atomization in response to patient respiratory forces exceeding a defined threshold. The impingement shield is operably associated with an intake valve such that when a negative threshold pressure is attained within the nebulizer, such as provided by an inhalation force provided by a respiring patient, said intake valve is moved from a closed to an open state establishing inhalation flow through the nebulizer. By rendering the air flow through the nebulizer as contingent upon exceeding a minimum negative pressure it is now possible to constrain atomization of a medicament to the combined operational status of the patient inhaling (to actuate the impingement shield into a second “open” or non-occluding state) and of the patient attaining a defined level of force during inhalation (allowing for opportunity to control inhalation airflow to coincide with deeper pulmonary penetration of medicinal nebula), thus offering enhanced therapy regimes. This combined operational status is particularly noteworthy in that prior art devices will not produce aerosol until there floating baffle is drawn all the way down, which is achieved only upon reaching a minimum inhalation vacuum pressure, the result of which is that it is possible with prior art devices for a patient to breath at a low inhalation flow rate insufficient to draw the floating baffle all the way down so that aerosol is produced. The unique and novel design of the current invention disallows any flow of air through the nebulizer until such time that shield has been drawn down, thus providing greater assurance that aerosol is delivered with each inhalation and being suitable for a broader range of patients. Because the current invention is suitable for a broad range of flow rates, it can be fabricated into a form with a very low nebulization flow rate (e.g. 1.0 l/min) that will maximize its sensitivity to the inhalation effort of the patient. [0011] In a further embodiment of the present invention, the nebulizer having a respiration responsive impingement shield as presented herein may further include an electronic sensor which is responsive to the position and duration of the impingement shield being in said first and second positions. By using a simple and conventional logic circuit, it is possible to indicate to the patient if insufficient inhalation or exhalation periods of occurred. The same logic circuit can be used to indicate optimum therapy duration to the patient, possible error conditions, or provide estimated dosages by flagging of visual and/or auditory cues. The same logic circuit can be used to transmit the same resulting information through electrically conducting or wireless means to a remote location for use as a clinical evaluation tool or to provide greater management of a patient's condition. DESCRIPTION OF THE SEVERAL VIEWS OF THE DRAWING(S) [0012] The invention will be more easily understood by a detailed explanation of the invention including drawings. Accordingly, drawings, which are particularly suited for explaining the inventions, are attached herewith; however, it should be understood that such drawings are for descriptive purposes only and as thus are not necessarily to scale beyond the measurements provided. The drawings are briefly described as follows: [0013] FIG. 1 is an exterior perspective view of an inhalation-controlled nebulizer in accordance with the present invention, wherein cross-sectional planes 3 - 3 and 4 - 4 are defined. [0014] FIG. 2 is an exploded perspective view of an inhalation controlled nebulizer as presented in FIG. 1 , wherein the impingement shield is presented. [0015] FIG. 3 is an exploded perspective view of an inhalation controlled nebulizer as presented in FIG. 1 , wherein the impingement shield is presented. [0016] FIG. 4 is a cross sectional side view of an inhalation controlled nebulizer with impingement shield in a non-aerosol producing position taken along line 3 - 3 . [0017] FIG. 5 is a cross sectional side view of an inhalation controlled nebulizer with impingement shield in an aerosol producing position taken along line 3 - 3 . [0018] FIG. 6 is a cross sectional side view of an inhalation controlled nebulizer with impingement shield in a non-aerosol producing position taken along line 4 - 4 . [0019] FIG. 7 is a left side view of an inhalation controlled nebulizer with impingement shield. [0020] FIG. 8 is a top end view of an inhalation controlled nebulizer with impingement shield. [0021] FIG. 9 is a perspective view of shield assembly with impingement shield. [0022] FIG. 10 is a sectional view of an alternate nozzle design applicable for use in an inhalation controlled nebulizer with impingement shield. [0023] FIG. 11 is a sectional view of an alternate nozzle design, hemispherical target surface, and impingement shield in a non-aerosol producing position applicable for use in an inhalation controlled nebulizer with impingement shield. [0024] FIG. 12 is a sectional view of an alternate nozzle design, flat target surface, and impingement shield in a non-aerosol producing position applicable for use in an inhalation controlled nebulizer with impingement shield. [0025] FIG. 13 is a perspective view of a patient tee assembly equipped with nebulizer tee, mouthpiece and exhalation filter. [0026] FIG. 14 is an exploded perspective view of a patient tee assembly equipped with nebulizer tee, mouthpiece and exhalation filter. [0027] FIG. 15 is a sectional view of a patient tee assembly equipped with nebulizer tee, mouthpiece and exhalation filter. [0028] FIG. 16 is a perspective view of an inhalation controlled nebulizer fitted with a patient tee assembly. LIST OF REFERENCE NUMERALS [0000] 4 Nebulizer Unit 6 Patient Tee Assembly 8 Insert 10 Upper Chamber 12 Lower Chamber 14 Liquid Reservoir 16 Aerosol Chamber 18 Internal Gas Conduit 20 Gas Inlet Port 22 Aerosol Outlet Port 23 Nebulizer Outlet Body 24 Secondary Shroud 26 Entrainment Orifice 28 Liquid Transfer Channels 30 Jet Orifice 31 Shield Assembly 32 Impingement Shield 33 Shield Yoke 34 Shield Flag 36 Yoke Mounting Flange 40 Biasing Support 42 Target Surface 44 Air Inlet Portal 46 Retention Ring 48 Liquid Backflow Guard 50 Nozzle Yoke 52 Ambient Chamber 54 Cylindrical Guide 56 Shield Cam 58 Shield Flow Ports 60 Adjustment Knob 62 Shield Minimum Flow Channel 64 Shield Flow Diverter 66 Adjustment Knob Actuating Teeth 68 Intermittent Ambient Gas Passage 70 Post Nebulization Filter 72 Mouthpiece 74 Nebulizer Tee 76 Check Valve Body 78 Check Valve Flapper 80 Filter Membrane 82 Mouthpiece Conduit 84 Mouthpiece Port 86 Nebulizer Port 88 Anti-Drool Chimney 90 Check Valve Port 92 Check Valve Flow Conduits 94 Flapper Retention Boss 96 Flapper Retention Orifice DETAILED DESCRIPTION OF THE INVENTION [0078] While the present invention is susceptible of embodiment in various forms, there is shown in the drawings and will hereinafter be described a presently preferred embodiment of the invention, with the understanding that the present disclosure is to be considered as an exemplification of the invention, and is not intended to limit the invention to the specific embodiment illustrated. [0079] Referring more specifically to the figures, for illustrative purposes the present invention is embodied in the apparatus generally shown in FIG. 1 through FIG. 16 . [0080] In FIG. 1 through 8 , therein is depicted nebulizer unit 4 . Nebulizer unit 4 is comprised of an upper chamber 10 and a lower chamber 12 . Upper chamber 10 has therein an insert 8 for allowing ambient air to be drawn through aerosol chamber 16 , chamber 10 and an aerosol outlet port 22 having a liquid backflow guard 48 , which is in fluid communication with a patient through a suitable mouth piece, mask, endotracheal tube, or patient tee assembly 6 as shown if FIG. 13 through 16 . Lower chamber 12 has therein a liquid reservoir 14 and a gas inlet port 20 . Gas inlet port 20 extends from an area exterior to lower chamber 12 whereby it is attached to a pressurized gas supply (not shown) and passes through liquid reservoir 14 , into secondary shroud 24 . In the embodiment shown, upper chamber 10 and lower chamber 12 are releasably affixed to one another so that liquid medicament can be introduced into liquid reservoir 14 . It is within the purview of the present invention that a liquid addition portal can be provided for introduction of liquid medicament, and in such case, upper chamber and lower chamber may be permanently affixed at the time of manufacture. [0081] Turning to FIGS. 4 , 5 , and 6 , therein is depicted secondary shroud 24 comprising a liquid transfer conduit 28 , jet orifice 30 , internal gas conduit 18 , and an orientation and construction support member nozzle yoke 50 . Gas provided to gas inlet port 20 is caused to pass through internal gas conduit 18 and onto jet orifice 30 , where with sufficient gas pressure (greater than 8 psig) provided to gas inlet port 20 , the gas jet emanating from jet orifice 30 will be a significant percentage of or be equal to the speed of sound. Liquid transfer conduit 28 forms an interstitial space between the external geometry of internal gas conduit 18 and the internal geometry of secondary shroud 24 and necessarily provides a free flowing path for liquid from the bottom of liquid reservoir 14 and entrainment orifice 26 . Secondary shroud 24 utilizes pressurized gas from gas inlet port 20 ejected through jet orifice 30 onto target surface 42 , wherein the proximity of surface 42 is sufficient to redirect the flow of the impinging gas radially thereby causing a vacuum across a proximal opening of entrainment orifice 26 , thus drawing liquid medicament from liquid reservoir 14 through liquid transfer conduit 28 and to entrain that liquid medicament into a continuous high velocity atomized radial fan. Although other nozzle configurations are also possible with the present invention, as exemplified in FIGS. 10 , 11 , and 12 , the fan shaped spray produced by the described nozzle has been found to be desirable. By way of the fluidic jet stream impacting upon target surface 42 and through the combined actions of jet dispersion, high jet momentum, and shear forces acting on introduced fluid, micro-droplets of liquid medicament are formed so long as there is liquid medicament to be entrained and a supply of gas through gas inlet port 20 . Target surface 42 may have a simple geometric or radiused cross sectional profile as well as compound combinations of differing geometric and/or radiused cross sectional profiles. In a preferred embodiment, target surface is of a flat, convex or hemispherical cross sectional profile. [0082] Positioned proximal to the target surface 42 is impingement shield 32 . Impingement shield 32 is capable of at least partially occluding the fluidic communication pathway between a nebulization area defined as the region between entrainment orifice 26 and fixed target surface 42 , and aerosol chamber 16 that is internal to insert 8 . When impingement shield 32 is in first position, the fluidic communication pathway between the nebulization area and aerosol chamber 16 is at least partially occluded. When the fluidic communication pathway is at least partially occluded by impingement shield 32 , micro-droplets produced by the interaction of entrainment orifice 26 with target surface 42 are slowed and the majority is caused to impact upon the impingement shield 32 . As micro-droplets of medicament are slowed and may impact upon the impingement shield 32 , the micro-droplets coalesce into macro-droplets, which in turn are un-respirable and return under gravity to liquid reservoir 14 . Impingement shield 32 can also be translated to a second position, wherein the fluidic communication pathway between the nebulization area and aerosol chamber 16 is not occluded, thus allowing micro-droplets of medicament to be released into aerosol chamber 16 . As a patient breathes, the impingement shield 32 oscillates between the first and second position. Although the geometry needed for impingement shield 32 to be effective at obstructing the flow of aerosol from the nebulization area to aerosol chamber 16 when obstruction is desired varies, it is preferred that impingement shield 32 be of sufficient height and placement in the obstructing position that the direction of travel of gas originating from jet orifice 30 , redirected by target surface 42 , and emanating from the nebulizer area be caused to change as a result of impingement shield 32 position prior to entering aerosol chamber 16 . It is furthermore preferred that the geometry around the nebulization area when impingement shield 32 is in an obstructing position redirect the gas stream preferentially upwards against the flow of gravity. In such manner, coalesced liquid accumulates to a greater degree within impingement shield 32 and around nebulization area, providing greater efficiency at the capture of aerosol particles leaving nebulization area and thereby causing a greater return of liquid medication to the reservoir. [0083] In an alternative design, secondary shroud 24 may comprise entrainment orifice 26 , liquid transfer conduit 28 , and jet orifice 30 ( FIG. 10 ). In accordance with the ejection nozzles taught by Lester in the aforementioned and incorporated patents of reference, as pressurized gas issues from jet orifice 30 into entrainment orifice 26 , liquid from reservoir 14 is drawn up liquid transfer conduit 28 . As liquid is drawn through liquid transfer conduit 28 , it passes through a flow control point (typically about 0.010 inch in height) and then into direct contact with, and becomes entrained within, the gas issuing from jet orifice 30 and is forcibly ejected from entrainment orifice 26 as a focused continuous stream of liquid entrained gas. The high velocity liquid entrained jet is oriented such that an optimized focal point is continuously achieved at a defined distance from the entrainment orifice 26 . As the liquid jet stream comes to the optimization point of the jet, the jet strikes a fixed target surface 42 . The liquid jet stream impacting upon target surface 42 , and through the combined actions of minimal jet dispersion and high jet momentum forms micro-droplets of liquid medicament so long as there is liquid medicament to be entrained and a supply of gas through gas inlet port 20 . It is within the purview of the present invention that one or more liquid entrained gas jets may be formed by secondary shroud 24 . The nozzle configuration of FIG. 10 is suitable with an array of different target surface 42 geometries including the hemispherical design shown in FIG. 11 , and the flat disc design shown in FIG. 12 . [0084] Within a central region of upper chamber 10 , there extends downwardly insert 8 . Insert 8 may be either an element integral to upper chamber 10 or separate element affixed to a central void within upper chamber 10 . In a preferred embodiment, insert 8 is generally round in cross section taken at a point parallel to a point of junction with lower chamber 12 . The insert 8 extends into a central void of upper chamber 10 and has a distal point that is proximal to secondary shroud 24 . At the distal point of insert 8 therein is an optional retention ring 46 that acts upon shield assembly 31 to maintain durable attachment to a biasing support 40 . [0085] Shield assembly 31 comprises impingement shield 32 shield yoke 33 , shield flag 34 and yoke mounting flange 36 ( FIG. 11 ). Impingement shield 32 is designed to at least partially prevent the movement of medicament nebula resulting from the interaction of the secondary shroud 24 and target surface 42 to aerosol chamber 16 when in at least one position, and to not prevent the movement of medicament nebula resulting form the interaction of secondary shroud 24 and target surface 42 to aerosol chamber 16 when in at least one other position. As depicted in the associated figures, a preferred embodiment of the impingement shield 32 is as a cylindrical ring having a height sufficient to at least partially occlude the fluid communication pathway and an interior diameter sufficient to circumscribe the outer diameter of a cylindrical cross-section region defined by the secondary shroud 24 and target surface 42 . While the secondary shroud 24 target surface 42 , and impingement shield 32 are depicted with circular cross-section, alternate cross-sectional geometries are possible so long as the impingement shield 32 can circumscribe the secondary shroud 24 /target surface 42 and at least partially occlude the associated fluidic communication pathway with aerosol chamber 16 . [0086] Attached to impingement shield 32 is shield yoke 33 . Shield yoke 33 connects the impingement shield 32 to a biasing support 40 , and maintains the impingement shield 32 in proper orientation relative to the secondary shroud 24 and target surface 42 . Shield yoke 33 may be connected to impingement shield 32 at one or more points and may be either a separate component durably affixed to impingement shield 32 or may be integrally formed with impingement shield 32 . Shield yoke 33 terminates at yoke mounting flange 36 . [0087] Yoke mounting flange 36 is connected to shield yoke 33 at one or more points and may be either a separate component durably affixed to shield yoke 33 or may be integrally formed with shield yoke 33 . Shield yoke may optionally include a shield flag 34 . Shield flag 34 extends outside upper chamber 10 , affording additional maintained orientation of the impingement shield 32 during operation of the nebulizer. In addition, shield flag 34 may be used to visually indicate operation of the nebulizer, or in the alternative, to trigger a simple and conventional logic circuit to electronically track operation of the nebulizer. [0088] Within upper chamber 10 , positioned in fluidic communication with aerosol chamber 16 , is a durably affixed biasing support 40 . Biasing support 40 is acted upon by respiration forces from the patient, wherein the force is translated into movement of the shield assembly 31 . Suitable biasing support 40 includes membranes which are responsive to changes in force or flow of air through the nebulizer, and include elastomeric materials such silicone, natural rubbers and blocked AB polymers. Further, biasing support 40 may be homogenous in construction, or comprised of two or more differing materials, having regions of same or dissimilar cross-sectional profiles, and same or differing extension, recovery and related physical performance properties. Biasing support 40 may further include one or more biasing members, such as coil or leaf spring, to further act upon the shield assembly 31 . [0089] Aerosol chamber 16 consists of the space immediately around the aerosol producing region defined generally as the region between and including target surface 42 and the face of secondary shroud 24 coincident with the exit plane of entrainment orifice 26 . Although in the preferred embodiment aerosol chamber 16 is encompassed by the internal geometry of insert 8 , the invention need not be limited to said configuration and other embodiments in which the outer limits of aerosol chamber 16 are defined by the internal geometry of upper chamber 10 and/or lower chamber 12 are possible without departing from the invention. [0090] At a point in upper chamber 10 , proximal to biasing support 40 and on a side opposite to biasing support 40 that is in continuous fluid communication with aerosol chamber 16 is air inlet portal 44 . Air inlet portal 44 provides fluid communication between the ambient environment and the interior of upper chamber 10 . Ambient chamber 52 consists of the volume of space that is in un-interrupted flow communication of air inlet portal 44 and is in interrupted flow communication with aerosol chamber 16 . Said interruption of flow communication between ambient chamber 52 and aerosol chamber 16 is caused by position of biasing support 40 such that during patient exhalation ambient chamber 52 and aerosol chamber are not in fluid communication, and fluid communication between ambient chamber 52 and aerosol chamber 16 only occurs in such instance that biasing support 40 has moved sufficient distance, either through force of inhalation or manual actuation, to allow impingement shield 32 sufficient movement so as to allow the movement of medicament nebula resulting from the interaction of the secondary shroud 24 and target surface 42 to aerosol chamber 16 . Thus a useful feature of the invention is that the minimum amount of air needed to be drawn by the patient for the impingement shield 32 to be in a non-occluding position is only the gas flow caused to flow through jet orifice 30 since no ambient air may be drawn through the nebulizer until such time that ambient chamber 52 and aerosol chamber 16 are in fluid communication In a preferred embodiment, upper chamber 10 includes cylindrical guide 54 , shield assembly 31 includes shield cam 56 , shield flow ports 58 , shield minimum flow channel 62 , and shield flow diverter 64 , and adjustment knob 60 includes adjustment knob actuation teeth 66 . Cylindrical guide 54 extends axially and centrally into the internal space of upper chamber 10 such that it encompasses shield cam 56 and shield minimum flow channel 62 . Adjustment knob actuation teeth 66 of adjustment knob 60 engage with shield cam 56 of shield assembly 31 such that rotation of adjustment knob 60 to the breath actuation mode allows for the travel of shield assembly up and down vertically with exhalation and inhalation of patient as herein described. Alternatively adjustment knob 60 may be rotated to the continuous mode, causing a different engagement of adjustment knob actuation teeth 66 with shield cam 56 such that shield assembly 31 is restricted to a down position so as to allow the movement of medicament nebula resulting from the interaction of the secondary shroud 24 and target surface 42 to aerosol chamber 16 regardless if the patient is inhaling or exhaling. Shield minimum flow channel 62 is located at the lowest point of shield cam 56 . Upon initiation of patient inhalation ambient air is not allowed to pass from ambient chamber 52 to aerosol chamber 16 due to the impediment of fluid communication caused by the position of shield minimum flow channel 62 with respect to cylindrical guide 54 . Upon patient inhalation becoming developed sufficiently to over-draw compressed gas provided through jet orifice 30 , shield assembly 31 will travel downwards allowing movement of medicament nebula resulting from the interaction of secondary shroud 24 and target surface 42 to aerosol chamber 16 , and causing shield minimum flow channel 62 to also travel downwards sufficiently to clear cylindrical guide 54 so as to create a gap between shield minimum flow channel 62 and cylindrical guide 54 and thereby forming intermittent ambient gas passage 68 and thus allowing the travel of ambient air through air inlet portal 44 , ambient chamber 52 , intermittent ambient gas passage 68 , shield flow ports 58 and aerosol chamber 16 . Upon patient exhalation, biasing support 40 is already in a position such that shield minimum flow channel 62 is in a position in relation to cylindrical guide 54 such that intermittent ambient gas passage 68 is not formed and thus exhaled gas is not allowed to escape out of or through nebulizer unit 4 . Furthermore, with greater exhalation effort biasing support 40 is pushed with greater force upon cylindrical guide 54 creating a greater seal and impediment to exhaled flow. Thus an optimum embodiment of the invention includes the use of a mouthpiece, mask, or endotracheal tube unit or assembly that is equipped with a route for the passage of exhalation gases, and more optimally equipped with a route for the passage of exhalation gases that is biased in favor of exhalation gases flowing from the patient to the ambient environment during exhalation and biased against the flow of ambient air to the patient during inhalation since this gas may be more readily and effectively provided through nebulizer unit 4 . An even more optimum configuration would include a filter through which exhaled gas was caused to pass through thus capturing exhaled particles not desired in the ambient environment. Patient tee assembly 6 shown in FIGS. 13 through 16 is one such optimum embodiment and is hereafter described in detail. Those skilled in the art can appreciate that a number of other configurations are possible that achieve the same objective of patient tee assembly 6 , none of which depart from the present invention. [0091] Further, without being constrained to specific theory, it is believed and understood by those skilled in the art that ambient air drawn through aerosol chamber 16 during inhalation allows for evaporation and reduction of size of droplets created by the aerosol producing region during patient inhalation, thus increasing the number of micro-droplets formed in the respirable range (i.e. 0-10 microns). The formation of copious of amounts of micro-droplets in the respirable range thereby forming an aerosol and filling out the remaining internal geometry of the invention and being drawn out by patient inhalation through aerosol outlet port 22 that is in fluid communication to the patient through use of a mouthpiece, mask, endotracheal tube or patient tee assembly 6 . [0092] If the inhalation triggered performance of nebulizer unit 4 is not desired, it is possible to override manually the impingement shield 32 through manual force applied to shield assembly 31 , such as by applying downward force to optional shield flag 34 . Force applied to shield flag 34 causes impingement shield 32 to move to the second, non-occluding position. [0093] In general practice with the nebulizer unit 4 in accordance with the present invention, supplied gas to inlet port 20 at a pressure of at least 8 psig at a flow rate of between 1 and 15 liters of gas per minute, with the range of 5 to 12 liters per minute inclusively being preferred and the range of 8 to 11 liters per minute inclusively being most preferred. The gas issues through a jet orifice 30 having a diameter in the range of 0.011 and 0.030 inches. One or more liquid transfer conduits 28 are provided in secondary shroud 24 so that a volume of liquid medicament can be provided for aerosolization. The cross sectional flow area through which entrained liquid flows prior to entering aerosol producing area being 2 to 12 times greater than the cross sectional flow area of jet orifice 30 . Impingement shield 32 having height at least as great as the distance from the exit plane of jet orifice 30 to the nearest point of target surface 42 . Impingement shield 32 having an inside perimeter such that the minimum cross sectional area for the flow of gas from the nebulization area to the aerosol chamber 16 when impingement shield 32 is in the obstructing position has an equivalent diameter that is less than twenty times the equivalent diameter of jet orifice 30 . [0094] It is within the purview of the present invention that an inhalation actuated nebulizer with impingement shield may be combined with one or more ancillary devices to further enhance respiratory therapy. A particularly advantageous embodiment includes use of a post nebulization filter, which upon a non-inhalation event, significantly reduces the release of residual medicament nebula from the nebulizer and patient. Such a post nebulization filter may operate by various modalities, including, but not limited to, size exclusion, impact, tortuous path, and depth filtration. Further, the post nebulization filter may include mechanically responsive means for cycling filter performance based on respiratory forces and recycling functions for returning captured medicament to the nebulizer liquid reservoir for reuse. Such combined nebulizer and post nebulizer filter(s) may be employed in situations where in the release of the medicament to the immediate atmosphere or ambient environment is expensive, deleterious to others, or of a controlled nature (i.e. palliative narcotics). One embodiment including said such post nebulizer filter configuration is shown in FIGS. 13-16 and is generally indicated by patient tee assembly 6 . Post nebulization filter 70 is optional for patient tee assembly 6 , which is an advantage of the demonstrated configuration of patient tee assembly 6 due to the additional expense of post nebulization filter 70 that is not needed in all instances. In addition to post nebulization filter 70 , patient tee assembly 6 also consists of mouthpiece 72 , nebulizer tee 74 , check valve body 76 , and check valve flapper 78 . Post nebulization filter 70 includes filter membrane 80 that is positioned such that all gas passing through the body of post nebulization filter 70 is caused to pass through filter membrane 80 . As known by those skilled in the art, filter membrane 80 may consist of a wide array of different materials depending on the expected need, including but not limited to glass fibers, cellulose acetate, cellulose nitrate, porous nylon and/or Teflon. Mouthpiece 72 includes a distal end shaped to comfortably fit in the patient's mouth and is equipped with a centrally positioned mouthpiece conduit 82 that allows gas to pass freely from either distal end to other. The distal end of mouthpiece 72 opposite the patient side is equipped with a tapered outside diameter, usually of 22 mm nominal dimension, and allows for press fit into mouthpiece port 84 of nebulizer tee 74 . Nebulizer tee 74 also consists of nebulizer port 86 , anti-drool chimney 88 , and check valve port 90 . Anti-drool chimney 88 is an internal feature to nebulizer tee 74 that is generally axially aligned with nebulizer port 86 but extends from the inside wall of nebulizer tee 74 towards the primary central axis sufficient distance that drool or bodily fluids excreted from the patients mouth that travel through mouthpiece conduit 82 are prevented from passing through nebulizer port 86 and into nebulizer unit 4 where they may be aerosolized and introduced to the patients lungs thus compromising the respiratory health of the patient. Nebulizer port 86 is a tapered diameter sized to engage with outer diameter of nebulizer outlet body 23 . Check valve body 76 includes check valve flow conduits 92 , and flapper retention boss 94 . Check valve flapper 78 is made of an elastic material, such as silicone, and includes flapper retention orifice 96 . When check valve flapper 78 is engaged with check valve body 76 by stretching check valve flapper 78 sufficiently for flapper retention orifice 96 to fit over flapper retention boss 94 the result is that check valve flapper 78 is held into place onto check valve body 76 so as to cover check valve flow conduits 92 . Upon insertion of the resulting check valve assembly into check valve port 90 , the result is a flow conduit that easily allows gas to pass out of nebulizer tee 74 through check valve port 90 and check valve flow conduits 92 , but largely prevents the entrainment of gas in the opposite direction, thus when engaged with nebulizer unit 4 , as shown in FIG. 16 , exhaled gas is preferentially directed out of nebulizer tee 74 through check valve port 90 . As herein described, nebulizer unit 4 is designed so that there is no route for exhaled gases to escape out of nebulizer unit 4 , but that inhaled gas is preferentially drawn through nebulizer unit 4 as previously described and not check valve port 90 . Upon placement of post nebulization filter 70 over check valve port 90 , all exhaled gas is thereby caused to pass through post nebulizer filter 70 whereby undesired particles are captured on filter membrane 80 prior to exhaled gas being released to the ambient environment, thereby preventing contamination of the ambient environment with potentially undesirable aerosol which can lead to undesired exposure of additional people present at the time of treatment. [0095] The general construction of functional elements of nebulizer unit 4 and patient tee assembly 6 , includes thermoset and thermoplastic polymers as well as alloys and blends within those plastic families. Additional performance and aesthetic modifying chemistries can be incorporated during manufacture or after component or device fabrication. Of particular interest, polymers having specific surface energies can be used in different aspects of nebulizer unit 4 depending upon the degree of liquid medicament wet-out is desired. The nebulizer unit 4 and patient tee assembly 6 of the present invention is not constrained by the mode of manufacture and may include known or developed methods in molding and machining technology. [0096] “Macro-droplets” are defined herein as being an individual unit of liquid medicament having an average diameter of greater than 10.0 micrometers and representing the predominant form of liquid medicament to pass through the aerosolization region and return to the reservoir. “Micro-droplets” are defined herein as being an individual unit of liquid medicament having an average diameter of less than or equal to 10.0 micrometers and being the predominant fraction of liquid medicament to pass through the aerosolization region and leave the nebulizer. [0097] “Equivalent diameter” for any one or combination of cross sectional areas or conduits, of any shape, is defined herein as being the square root of the product of the cross sectional area, or sum of cross sectional areas for more than one cross sectional area or conduit, and four divided by Pi. [0098] For the purposes of general background in aerosol technology and as indicia for the level of understanding resident in one skilled in the art of aerosol technology, the following references are incorporated by reference in their entireties as nonessential matter: “Aerosol Technology”, Hinds, 1982 by Wiley and Sons, ISBN 0-471-087726-2; “Inhalation Aerosols”, First Ed. Hickey, 1996 by Informa Healthcare, ISBN 0-8247-9702-7; and “The Mechanics of Aerosols”, Fuchs et al., 1989 by Dover Publications, ISBN 0-486-66055-9. Example [0099] A first embodiment nebulizer device in accordance with the present invention was fabricated and tested. [0100] Device Dimensions: [0000] Nebulizer Unit height: 3.75 inches Nebulizer Unit width/diameter: 1.88 inches Shield Assembly Travel Distance: 0.17 inches Inhalation port dimensions: 22 mm ISO ID “Respiratory Conduction Mouthpiece” Liquid reservoir maximum volume fill: 6.0 milliliters Aerosol Outlet Port Diameter: 0.90 inches [0101] In practice with the above example of a preferred embodiment, nebulizer unit 4 is supplied gas to inlet port 20 at a pressure of at least 8 psig, and more preferably at least 13 psig, at a flow rate of between 1 and 15 liters of gas per minute, with the range of 5 to 12 liters per minute inclusively being preferred and the range of 8 to 11 liters per minute inclusively being most preferred. The gas issues through a jet orifice having a diameter in the range of 0.011 to 0.030 inches, and more preferably in the range of 0.019 to 0.026 inches. The ratio of the cross sectional area that the liquid flows from into the aerosol producing region versus the cross sectional area of the jet orifice being preferably between 2 and 12 and more preferably between 4 and 8. The ratio of the distance from the exit plane of jet orifice to target surface and the diameter of entrainment orifice being preferably in the range of 0.1 and 0.8, and being more preferably in the range of 0.4 and 0.6. The ratio of the diameter of target surface and the entrainment orifice being preferably at least 1.0 and being more preferably at least 1.4. The height of impingement shield being preferably at least as large as the distance from the exit plane of the jet orifice to the nearest point of the target surface and being so positioned during exhalation or non-use to cause a redirection gas emanating from the target surface and jet orifice. The equivalent inside diameter of the minimum cross sectional area through which gas needs to flow from the nebulization area to the aerosol chamber as a result of the position of the impingement shield during exhalation or non-use being preferably less than the twenty times the equivalent diameter of the jet orifice, and being more preferably less than ten times the equivalent diameter of the jet orifice. When operating in the above format, it is possible to complete a 3 ml dosage of liquid medicament in less than four minutes of inhalation time at a nominal flow rate of 8 to 10 liters per minute under conditions in which aerosol is only delivered to the patient during inhalation, and by which viscous or thick medications can be delivered efficiently. Testing [0102] Inhalation Effort Evaluation Protocol a. Obtain a Competitive Commercial Model w/ Movable Baffle Technology (e.g. AeroEclipse Nebulizer by Trudell Medical International) and a representative example of the present invention as disclosed herein. b. Place 3 ml of water into the nebulizer. c. Connect to a Harvard pump (Model #). d. Run at specified flow rate of 8 l/min. e. Keep the inhalation (“I”) to exhalation (“E”) ratio at 1:1. f. Adjust the tidal volume to 250 ml. g. Start at a respiratory rate of 20 breaths per minute, and adjust the respiratory rate down until the AeroEclipse fails to produce aerosol during inhalation due to insufficient inhalation pressure. h. Replace the AeroEclipse with an example in accordance with the present invention that incorporates an inhalation optimization feature as described. i. Note whether any aerosol is being produced at the same settings and repeat the test. [0112] Inhalation Effort Results [0113] At a tidal volume of 250 ml, and an I:E ratio of 1:1, the AeroEclipse nebulizer stopped producing aerosol at a breathing rate of 12 breaths per minute (bpm) or less. Looking into the nebulizer it was visually determined that the moving critical component was not closing sufficient distance to allow aerosolization. The nebulizer in accordance with the present invention, including inhalation optimization feature, produced aerosol at 12 bpm, and continued to produce aerosol during inhalation at the same tidal volume and I:E ratio down to a breathing rate of 5 bpm. [0114] Surface Tension Evaluation Protocol a. Obtain a Competitive Commercial Model w/ Movable Baffle Technology (e.g. [0116] AeroEclipse Nebulizer by Trudell Medical International) and a representative example of the present invention as disclosed herein. b. Mix albumin protein into water at 120 mg/ml. c. Place 5 ml of protein mixture into AeroEclipse. d. Connect nebulizer to a simulated patient condition of TV=425 ml, Respiratory Rate=16 bpm, and I:E=1:1. e. Note if aerosol stops production during exhalation. f. Repeat using an example in accordance with the present invention. [0122] Surface Tension Results [0123] When filled with 5 ml of albumin protein (120 mg/ml) the AeroEclipse nebulizer produced clearly visible aerosol during both inhalation and exhalation (indicative of breath-actuation mode failure). The nebulizer in accordance with the present invention, filled with 5 ml of albumin protein (120 mg/ml), cycled between aerosol production during inhalation and a proper cessation in aerosol production during simulated exhalation. [0124] Adhesion Performance Degradation Evaluation Protocol a. Obtain a Competitive Commercial Model w/ Movable Baffle Technology (e.g. AeroEclipse Nebulizer by Trudell Medical International) and a representative example of the present invention as disclosed herein. b. Mix acetylcysteine at 10 mg/ml in water (same concentration as recommended for the brand name equivalent Mucomyst) c. Place 8 ml of acetylcysteine in AeroEclipse nebulizer. d. Connect nebulizer to a simulated patient condition of TV=425 ml, Respiratory Rate=16 bpm, and I:E=1:1. e. Using an aerosol trap on an inhalation limb to capture aerosolized medication delivered during inhalation. f. Obtain an initial weight on nebulizer to determine gross output over the first minute. g. Perform prescribed performance test for the first minute. h. Run the nebulizer for 8 minutes. i. Repeat the performance tests for another minute. j. Compare the ratio of inhalation to gross output for the first minute to the last minute. k. Repeat with example in accordance with the present invention. I. Repeat for both devices using distilled water as a control test. [0137] Adhesion Performance Degradation Results [0138] When filled with 8 ml of Acetylcysteine 10 mg/ml (Mucomyst) and run over a total time of 8 minutes, the AeroEclipse showed a 25% reduction in the amount of medication delivered in the last minute as compared to the amount of medication delivered in the first minute. The shield nebulizer equivalent showed no significant drop in performance under the same conditions and during the same time period. Neither device showed any reduction in output when testing was repeated using distilled water. [0139] Although the description above contains many details, these should not be construed as limiting the scope of the invention but as merely providing illustrations of some of the presently preferred embodiments of this invention. Therefore, it will be appreciated that the scope of the present invention fully encompasses other embodiments, which may become obvious to those skilled in the art. In the appended claims, reference to an element in the singular is not intended to mean “one and only one” unless explicitly so stated, but rather “one or more.” All structural, chemical, and functional equivalents to the elements of the above-described preferred embodiment that are known to those of ordinary skill in the art are expressly incorporated herein by reference and are intended to be encompassed by the disclosure and present claims. Moreover, it is not necessary for a device or method to address every problem sought to be solved by the present invention, for it to be encompassed by the disclosure and present claims. Furthermore, no element, component, or method step in the present disclosure is intended to be dedicated to the public regardless of whether the element, component, or method step is explicitly recited in the claims. No claim element herein is to be construed under the provisions of 35 U.S.C. 112, sixth paragraph, unless the element is expressly recited using the phrase “means for.”

The present invention is directed generally to a nebulizer for the formation of micro-droplets from liquid medicaments for respiratory patient treatment, and more specifically, to a baffled nebulizer wherein a static baffle used to form an atomized medicament is proximal to a shied which responds to patient respiration force to oscillate from an aerosol flow occluding position to an aerosol flow open position. During inhalation, the shield moves into a first registration format to allow passage of the atomized medicament (nebula) to the patient. During exhalation/non-use, a biasing pressure maintains said shield in a second registration format such that the nebula is retarded from passing to the patient and is coalesced into macro-droplets which return to a supply reservoir for re-atomization. The present nebulizer design is particularly adaptable for controlling atomization in response to patient respiratory forces exceeding a defined threshold; allowing for opportunity to control inhalation airflow and enhanced therapy regimes.

BACKGROUND OF THE INVENTION 1. Field of the Invention The present invention relates to a heat-resistant adhesive sheet for temporarily fixing chips used in a method for fabricating substrateless semiconductor packages that do not use a metal lead frame. 2. Description of the Related Art Among LSI packaging technologies, Chip Size/Scale Package (CSP) technologies have recently come into attention. Among those technologies, substrateless semiconductor package technology such as Wafer Level Package (WLP) is attractive in terms of packaging density and size reduction. In a WLP fabrication method, multiple semiconductor Si wafer chips orderly arranged without the use of a substrate are encapsulated with an encapsulation resin at a time and then the wafer is diced into individual structures. Thus the method enables packages smaller than conventional ones that use a substrate to be fabricated efficiently. In such a WLP fabrication method, chips, which are conventionally fixed on a substrate, need to be fixed on an alternative supporter. Furthermore, since the chips need to be unfixed after the chips have been encapsulated with resin and formed into individual packages, the supporter need to be removable, instead of permanent bonding fixation. Attention is therefore being given to an approach to using an adhesive sheet as such a supporter for temporarily fixing chips. For example, Japanese Patent Laid-Open No. 2001-308116 describes a chip electronic component fabrication method that includes the steps of: attaching acrylic resin adhesion means onto a substrate, the adhesive means being adhesive before processing but the adhesion strength decreases after the processing; fixing a plurality of semiconductor chips of the same type or different types onto the adhesion means with an electrode-formed surface down; coating a whole area including interspaces between the plurality of semiconductor chips of the same type or different types with a protective material; applying predetermined processing to reduce the adhesion strength of the adhesion means and peeling off a pseudo wafer on which the semiconductor chips are fixed from the semiconductor chips; and cutting the protective material between the plurality of semiconductor chips of the same type or different types to separate the semiconductor chips or chip electronic components. Japanese Patent Laid-Open No. 2001-313350 describes a hip electronic component fabrication method that includes the steps of: attaching acrylic resin adhesion means onto a substrate, the adhesive means being adhesive before processing but the adhesion strength decreases after the processing; fixing a plurality of semiconductor chips of the same type or different types onto the adhesion means with an electrode-formed surface down; coating a whole area including interspaces between the plurality of semiconductor chips of the same type or different types with a protective material; removing the protective material from the area from the side opposite of the electrode-formed side to the side opposite of the semiconductor chips; applying predetermined processing to reduce the adhesion strength of the adhesion means and peeling off a pseudo wafer on which the semiconductor chips are fixed from the semiconductor chips; and cutting the protective material between the plurality of semiconductor chips of the same type or different types to separate the semiconductor chips or chip electronic components. According to these methods, the protection of the chips also protects the chips during mounting/handling after dicing and the packaging density can be improved. Japanese Patent Laid-Open No. 2008-101183 describes a dicing/die bonding tape including an adhesive layer containing epoxy resin and acrylic rubber and a method for bonding a semiconductor device resulting from dicing onto a supporter. Obviously, the method is not intended for substrateless semiconductor devices and the adhesive layer is chosen by taking into consideration the adhesiveness to a substrate. The following problems can arise with the following method for fabricating a substrateless semiconductor package using an adhesive sheet as a temporary supporter, which do not arise with fabrication of semiconductor packages using a lead frame. The problems will be described below with reference to FIG. 1 , which illustrates the substrateless semiconductor device fabrication method. Multiple chips 1 are attached onto a heat-resistant adhesive sheet 2 for semiconductor device fabrication. The heat-resistant adhesive sheet 2 includes an adhesive layer 12 on one side and a substrate fixing bond layer 13 on the other side. The heat-resistant adhesive sheet 2 for semiconductor device fabrication is fixed on a substrate 3 to form a structure depicted in part (a) of FIG. 1 . Alternatively, the heat-resistant adhesive sheet 2 for semiconductor device fabrication is attached onto a substrate 3 and chips 1 are fixed on the heat-resistant adhesive sheet 2 to form the structure depicted in part (a) of FIG. 1 . Then, the chips 1 on the structure depicted in part (a) are encapsulated together with an encapsulation resin 4 to form a structure illustrated in part (b) of FIG. 1 . Then, as illustrated in part (c), the heat-resistant adhesive sheet 2 , together with the substrate 3 , is removed from the multiple chips 1 encapsulated with the encapsulation resin 4 , or the multiple chips 1 encapsulated with the encapsulation resin 4 and the heat-resistant adhesive sheet 2 are removed together from the substrate 3 and then the heat-resistant adhesive sheet 2 for semiconductor device fabrication is removed from the chips 1 , thereby obtaining the multiple chips 1 encapsulated with the encapsulation resin 4 . Electrodes 5 are formed in desired positions on surfaces of the chips 1 encapsulated with the encapsulation resin 4 that are exposed on the side on which the heat-resistant adhesive sheet 2 for semiconductor device fabrication is provided, thereby forming a structure depicted in part (d). For the step of dicing, a dicing tape 8 having a dicing ring 7 provided on its encapsulation resin 4 side as required is bonded to the structure to fix the chips 1 encapsulated with the encapsulation resin 4 . The structure is diced with a dicing blade 6 as depicted in part (e) to ultimately provide multiple substrateless packages in which multiple chips are encapsulated with the resin as depicted in part (f). During the resin encapsulation, the heat-resistant adhesive sheet 2 for semiconductor device fabrication illustrated in FIG. 2( a ) can be deformed in planar directions due to expansion and elasticity of a base material layer or the adhesive layer of the heat-resistant adhesive sheet 2 for semiconductor device fabrication as illustrated in FIG. 2( b ). Accordingly, the positions of the chips 1 provided on the heat-resistant adhesive sheet 2 for semiconductor device fabrication can move. As a result, when the electrodes are provided on the chips 1 , relative positional relationship between the chips 1 and the electrodes would have changed from the originally designed relationship. Furthermore, during dicing of the chips 1 encapsulated with resin, the dicing line based on the positions of the chips 1 predetermined for the dicing step would be different from the dicing line required by the actual positions of the chips 1 . Consequently, the positions of chips encapsulated in the packages resulting from dicing would vary from one package to another and a subsequent electrode formation step would not smoothly be performed and partially encapsulated packages would result. When the heat-resistant adhesive sheet 2 for semiconductor device fabrication is peeled away from the resin-encapsulated chips, an adhesive formed on the chip side of the heat-resistant adhesive sheet 2 for semiconductor device fabrication exhibits heavy peeling from the chips and the encapsulation resin, depending on the properties of the adhesive. Therefore it can be difficult to peel off the heat-resistant adhesive sheet 2 for semiconductor device fabrication, or adhesive deposits 9 as illustrated in FIG. 3 can occur or static electricity can build up during peeling. As peeling becomes difficult, more time is required accordingly. Heavy peeling therefore can lead to reduction in productivity. Adhesive deposits 9 can inhibit a subsequent step such as electrode formation. Static electricity build-up caused by peeling leads to a problem due to adhesion of dust in a subsequent step. As has been described, chips can be displaced from specified positions by pressure applied during resin encapsulation because the chips are not properly held in the substrateless semiconductor package fabrication method using a heat-resistant adhesive sheet 2 for semiconductor device fabrication as a supporter for temporary fixture. When the heat-resistant adhesive sheet 2 for semiconductor device fabrication is peeled off, packages can be damaged by adhesion strength to the chips increased by curing of the encapsulation material or heat. These problems are specific to substrateless semiconductor device fabrication methods not suffered by other methods such as the method described in Japanese Patent Laid-Open No. 2008-101183. SUMMARY OF THE INVENTION Means for solving the problems is as follows. A heat-resistant adhesive sheet for semiconductor device fabrication is attached to a substrateless semiconductor chip when the substrateless semiconductor chip is encapsulated with resin. The heat-resistant adhesive sheet includes a base material layer and an adhesive layer. When bonded, the adhesive layer has an adhesion strength to SUS304 greater than or equal to 0.5 N/20 mm. The adhesive layer cures by stimulation received in a period between completion of the bonding and completion of resin encapsulation so that the peel strength of the adhesive layer from package reduces to 2.0 N/20 mm or less. The adhesive layer may be a radiation curable layer that cures under radiation such as ultraviolet or may be a radiation curable layer that cures by heat. Furthermore, there is provided a semiconductor device fabrication method that resin-encapsulates a substrateless semiconductor chip using the heat-resistant adhesive sheet 2 for semiconductor device fabrication, instead of a metal lead frame. The present invention provides a heat-resistant adhesive sheet 2 for semiconductor device fabrication that is used in a substrateless semiconductor package fabrication method (such as WLP fabrication method) that does not use a metal lead frame, and a semiconductor device fabrication method that uses the sheet. According to the present invention, when substrateless semiconductor chips that do not use a metal lead frame are encapsulated with resin, the chips are held so that the chips are not displaced from specified positions. Adhesive deposit is not left after the heat-resistant adhesive sheet 2 for semiconductor device fabrication has been used. Therefore, gas is not generated and an adhesive does not melt and attach during heating. Consequently, interconnects can be reliably provided, leading to improvement of fabrication yield of the semiconductor package and reduction in contamination with adhesive deposits that would otherwise occur during peeling of the sheet. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is a schematic diagram of a method for fabricating a substrateless package; FIG. 2 is a diagram illustrating deformation of a heat-resistant adhesive sheet for semiconductor device fabrication on which chips are mounted by heat during encapsulation with encapsulation resin; FIG. 3 is a diagram illustrating static electricity build-up and adhesive deposits that occur when a heat-resistant adhesive sheet for semiconductor device fabrication is removed; and FIG. 4 is a cross-sectional view of a heat-resistant adhesive sheet for semiconductor device fabrication according to the present invention. Description of Symbols  1: Chip  2: Heat-resistant adhesive sheet for semiconductor device fabrication  3: Substrate  4: Encapsulation resin  5: Electrode  6: Dicing blade  7: Dicing ring  8: Dicing tape  9: Adhesive deposit 10: Flat, smooth peeling sheet 11: Base material layer 12: Adhesive layer 13: Substrate fixing bond layer DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS To achieve the objects described above, the present inventors have diligently studied materials and configurations for a heat-resistant adhesive sheet for semiconductor device fabrication. As a result, the inventors have found that the objects described above can be achieved by using a heat-resistant adhesive sheet for semiconductor device fabrication that includes an adhesive layer having a peel strength in certain ranges before and after curing and by inducing a curing reaction of the adhesive layer by heat history experienced by the adhesive sheet in the period from completion of bonding to completion of resin encapsulation, and have made the present invention. The present invention relates to a heat-resistant adhesive sheet for semiconductor device fabrication that is attached to a substrateless semiconductor chips that do not use a metal lead frame when the chips are encapsulated with resin. The heat-resistant adhesive sheet for semiconductor device fabrication includes a base material layer and an adhesive layer. The adhesive sheet has been designed so that the adhesive layer cures in a period from completion of bonding to completion of resin encapsulation. In a specific heat-resistant adhesive sheet for semiconductor device fabrication, the adhesive layer contains a radiation curable adhesive or an adhesive that cures by heat. The radiation curable adhesive is preferably an ultraviolet (UV) curable adhesive. The heat-resistant adhesive sheet for semiconductor device fabrication has been designed so that the adhesive layer cures in a period from completion of bonding to completion of resin encapsulation. Curing means is not limited to specific means. For example, a UV curable adhesive may be contained in part or all of the adhesive layer and a curing reaction of the adhesive layer containing the UV curable adhesive may be induced by heating during the resin encapsulation step. An embodiment of a heat-resistant adhesive sheet 2 for semiconductor device fabrication of the present invention will be described in detail with reference to FIG. 4 . FIG. 4 is a cross-sectional view of the heat-resistant adhesive sheet 2 for semiconductor device fabrication including a base material layer 11 and an adhesive layer 12 . A substrate fixing bond layer 13 can be formed on the surface of the base material layer on the side where the adhesive layer is not provided, so that the heat-resistant adhesive sheet 2 having chips 1 fixed on the adhesive layer 12 can be fixed on the substrate 3 . The adhesive layer 12 is made of an adhesive such as an acrylic-based adhesive, an adhesive containing rubber and epoxy resin components, or a silicone resin-based adhesive. Flat, smooth peeling sheets 10 that protect the surfaces of the adhesive layer 12 and the substrate fixing bond layer 13 can also be provided. The heat-resistant adhesive sheet 2 for semiconductor device fabrication of the present invention will be described below. [Adhesive Layer 12 ] The adhesive of the adhesive layer 12 may be any material that is heat-resistant. Specifically, the adhesive may be any of various adhesives such as an acrylic-based adhesive, an adhesive containing rubber and epoxy resin components, and a silicone resin-based adhesive. The adhesive layer 12 is capable of reliably fixing chips before resin encapsulation and allows the heat-resistant adhesive sheet 2 for semiconductor device fabrication to be peeled away from resin-encapsulated chips without leaving adhesive deposits or damage to the resin. To that end, the adhesive layer 12 needs to have an initial adhesion strength higher than or equal to 0.5 N/20 mm to SUS304 and a peel strength lower than or equal to 2.0 N/20 mm from package after curing. The adhesive layer 12 has an initial adhesion strength of preferably higher than or equal to 0.6 N/20 mm to SUS304, more preferably higher than or equal to 0.8 N/20 mm and has a peel strength preferably lower than or equal to 1.0 N/20 mm from package, more preferably lower than or equal to 0.85 N/20 mm, yet more preferably lower than or equal to 0.5 N/20 mm. The peel strength from package herein is primarily the peel strength from cured encapsulation resin of the package because the cured encapsulation resin has a greater area adhered to the adhesive layer 12 of the heat-resistant adhesive sheet 2 for semiconductor device fabrication than the area adhered to the chips, although the adhesive layer 12 of the heat-resistant adhesive sheet 2 for semiconductor device fabrication is adhered to both of the cured encapsulation resin and the chips. Acrylic Resin-Based Adhesive Examples of the acrylic resin-based adhesive include an acrylic copolymer derived from copolymerization of monomers containing at least alkyl(meth)acrylate. Examples of alkyl(meth)acrylate include methyl(meth)acrylate, ethyl(meth)acrylate, butyl(meth)acrylate, isoamyl(meth)acrylate, n-hexyl(meth)acrylate, 2-ethyl hexyl(meth)acrylate, isooctyl(meth)acrylate, isononyl(meth)acrylate, decyl(meth)acrylate, and dodecyl(meth)acrylate. Alkyl(meth)acrylate is an alkyl acrylate and/or alkyl methacrylate. The same applies to (meth) in all compound names herein. The acrylic resin-based adhesive may contain any appropriate known cross-linking agent. Examples of the cross-linking agent include an isocyanate cross-linking agent, an epoxy cross-linking agent, an aziridine-based cross-linking compound, and a chelate-based cross-linking agent. The content of the cross-linking agent is not limited. Specifically, for example the content of the cross-linking agent is preferably 0.1 to 15 pts.wt., more preferably 0.5 to 10 pts.wt, for 100 pts.wt. of the content of the acrylic-based copolymer stated above. If the content of cross-linking agent is greater than or equal to 0.1 pts.wt., the adhesive layer 12 will have an appropriate level of viscoelasticity and therefore the adhesion strength of the adhesive layer 12 to conductive patterns on the surfaces of the chips or to the encapsulation resin 4 will not become excessively high. Therefore the encapsulation resin 4 will not be peeled off or damaged when the heat-resistant adhesive sheet 2 for semiconductor device fabrication is peeled off. Furthermore, portions of the adhesive layer 12 cannot adhere to and remain on the conductor patterns on the surfaces of the chips or encapsulation resin 4 . On the other hand, if the content of the cross-linking agent is less than or equal to 15 pts.wt, cracks of the adhesive layer 12 due to excessive cure of the adhesive layer 12 cannot occur. By containing a radiation curable adhesive in the adhesive layer 12 , the resin can be cured by heat during resin encapsulation and/or post cure. The radiation curable adhesive is preferably a UV curable compound that can be cured under irradiation of ultraviolet and is efficiently reticulated three-dimensionally by UV irradiation. Examples of such a UV curable compound include trimethylolpropane triacrylate, tetramethylolmethane tetraacrylate, pentaerythritol monohydroxy pentaacrylate, 1,4-butylenes glycol diacrylate, 1,6-hexanediol diacrylate, and polyethylene glycol diacrylate. A UV curable compound that contains more functional groups having an unsaturated bond which is a reaction site is more preferable. Especially preferable is a compound containing six or more functional groups. These compounds can be used either singly or in combination. The UV curable compound may be a UV curable resin. Examples of the UV curable compound include a photosensitive reactive group-containing polymer or oligomer, such as ester(meth)acrylate, urethane(meth)acrylate, epoxy(meth)acrylate, melamine(meth)acrylate, and acrylic resin (meth)acrylate, thiol-ene added resin and cationic photopolymerization-type resin having an allyl group at the molecular end thereof, a cinnamoyl group-containing polymer such as polyvinyl cinnamate, diazotized amino novolac resin, and an acrylamide type polymer. Examples of polymer reactive to UV include epoxidized polybutadiene, unsaturated polyester, polyglycidyl methacrylate, polyacrylamide, and polyvinyl siloxane. The content of the UV curable compound is preferably 5 to 500 pts.wt., more preferably 15 to 300 pts.wt., and especially preferably 20 to 150 pts.wt., for 100 pts.wt. of adhesive. Rubber-Component and Epoxy Resin-Based Adhesive (Rubber Component) Examples of the rubber component of the rubber and epoxy resin components containing adhesive agent constituting the adhesive layer 12 include rubber components conventionally used in epoxy-based bonds, such as NBR (acrylonitrile-butadiene rubber), acrylic rubber, acid terminated nitrile rubber, and thermoplastic elastomer. Examples of commercially available rubber components include NiPol 1072 (from Zeon Corporation) and Nipol-AR51 (from Zeon Corporation). Among them, NBR is preferably used in terms of compatibility with epoxy resin. In particular, NBR having an acrylic nitrile content of 10 to 50% is preferable. The aim of addition of the rubber component is to provide flexibility to the adhesive. However, the heat resistance decreases as the content of the rubber component increases. In view of this, the proportion of the rubber component in the organic substance in the adhesive layer 12 is preferably 20 to 60 wt %, more preferably 30 to 50 wt %. A proportion of 20 wt % or more adds flexibility to the adhesive layer 12 and provides a good machinability in cutting of the adhesive sheet. A proportion of 60 wt % or less provides a sufficient heat resistance and can prevent adhesives deposit. (Epoxy Resin Component) Examples of the epoxy resin component used with the rubber component include a compound containing two or more epoxy groups in its molecule, such as glycidyl amine-type epoxy resin, bisphenol F-type epoxy resin, bisphenol A-type epoxy resin, phenol novolac-type epoxy resin, cresol novolac-type epoxy resin, biphenyl-type epoxy resin, naphthalene-type epoxy resin, aliphatic epoxy resin, acrylic epoxy resin, heterocyclic epoxy resin, spiro-ring-containing epoxy resin, and halogenated epoxy resin. These components can be used either singly or in combination. Among these, bisphenol A-type epoxy resin is especially preferable in terms of peelability from the encapsulation resin 4 after the encapsulation step. The proportion of the epoxy resin used is preferably 40 to 80 wt %, more preferably 50 to 70 wt %, for 100 pts.wt. of organic substance. A proportion of 40 wt % or more can provide a sufficient heat resistance to the adhesive layer 12 by curing and a proportion of 80 wt. or less can provide flexibility and improve the machinability. The epoxy resin has a weight per epoxy equivalent of less than or equal to 1000 g/eq, preferably less than or equal to 500 g/eq. A weight per epoxy equivalent of 1000 g/eq or less can prevent adhesive deposits during peeling after the encapsulation step because the crosslink density is not decreased and bonding strength after curing falls within an appropriate range. Silicone Resin-Based Adhesive Preferably, the silicone resin-based adhesive has a storage elastic modulus greater than or equal to 5.0×10 3 N/cm 2 at 200° C. The silicone resin-based adhesive layer 12 is preferably 1 to 50 μm thick and preferably has an adhesion strength in the range of 0.05 to 4.0 N/20 mm after heating at 200° C. in conformity with JIS C2107. Such an adhesive layer 12 may be an addition-polymerized silicone resin adhesive layer 12 cross-linked by an organo polysiloxane structure, preferably a dimethyl polysiloxane structure, and an unsaturated group such as a vinyl group, and a SiH group and cured by a platinum-based catalyst, or a silicone resin-based adhesive layer 12 obtained by curing by an organic peroxide such as BPO. However, an addition-polymerized silicone resin adhesive layer is preferable in terms of heat resistance. In that case, the crosslink density can be adjusted according to the density of the unsaturated groups by taking into account the adhesion strength that can be obtained. Heating or other processing is required for enabling addition polymerization to form the silicone resin-based adhesive layer 12 . Since a large difference between the adhesion strength to silicone and the adhesion strength to the encapsulation resin can cause static electricity build-up, it is desirable that the adhesion strengths to the silicone and the encapsulation resin be as close to each other as possible. Therefore the adhesion strengths to them need to be within the ranges enumerated above. Since the thermal expansivity of the silicone resin-based adhesive layer is small, the amount of displacement of chips after resin encapsulation is small. The amount is less than or equal to 3 mm, preferably 0.1 mm. If packages are contaminated with gas generated from the adhesive layer 12 of the heat-resistant adhesive sheet 2 for semiconductor device fabrication in the step of curing the encapsulation resin by heat, the reliability of the packages can degrade due to poor plating during rewiring. Therefore, a reduction in the weight of the adhesive layer 12 at 180° C. is maintained at less than or equal to 3.0 wt %, preferably less than or equal to 2.0 wt %. Other Components In addition to the components described above, appropriate additives may be added to the adhesive layer 12 , such as a UV polymerization initiator for curing a UV curable compound or a thermal polymerization initiator, as required. An appropriate known polymerization initiator can be chosen as the UV polymerization initiator. The content of the polymerization initiator is preferably 0.1 to 10 pts.wt., more preferably 1 to 5 pts.wt., for a 100 pts.wt. of adhesive. A UV polymerization accelerator may be added along with the UV polymerization initiator as required. Other optional components may be added, such as a plasticizer, a pigment, a dye, an anti-aging agent, an antistatic agent, and other additives, such as filler, for improving the physical properties, such as the elastic modulus of the adhesive layer 12 . In particular, addition of conductive filler can improve the elastic modulus of the adhesive layer 12 and also prevent static electricity build-up caused by peeling. In particular, an anti-aging agent may be added in order to prevent deterioration at high temperatures. The thickness of the adhesive layer 12 is normally on the order of 1 to 50 μm, more preferably 5 to 30 μm. [Base Material Layer 11 ] A material for the base material layer 11 is not limited to a particular type. Any base material that is heat-resistant under heating conditions during resin encapsulation can be used. Since the resin encapsulation step is performed typically at a temperature around 175° C., a base material used is preferably heat-resistant so that the base material does not significantly contract or the base material layer 11 itself is not damaged at such temperatures. Accordingly, the base material has preferably a linear thermal expansion coefficient of 0.8×10 −5 to 5.6×10 −5 /K at a temperature of 50 to 250° C. If a base material that has a glass transition temperature lower than the heating temperature for curing the encapsulation resin 4 is used as the base material, the linear thermal expansion coefficient of the base material in a range of temperatures higher than the glass transition temperature will be higher than the linear thermal expansion coefficient in a range of temperatures lower than the glass transition temperature. Accordingly, displacement of adhered chips 1 from specified positions will increase. In addition, a uniaxially- or biaxially-stretched base material, which was stretched at a temperature higher than its glass transition temperature, starts contracting at a temperature lower than the glass transition temperature, which also increases displacement from the specified positions of the adhered chips. The positional accuracy of the chips can be improved by choosing a material that has a glass transition temperature higher than 180° C. as the material of the base material layer 11 of the heat-resistant adhesive sheet 2 which is attached to the substrateless semiconductor chips without a metal lead frame when the chips are encapsulated with resin can. Examples of such a base material include heat-resistant plastic films such as a polyethylene naphthalate (PEN) film, polyethylene sulfone (PES) film, polyetherimide (PEI), a polyethersulfone (PSF) film, a polyphenylene sulfide (PPS) film, a polyether ether ketone (PEEK) film, a polyarylate (PAR) film, an aramid film, and liquid crystal polymer (LCP). If the temperature at which the resin encapsulation is performed is less than or equal to 150° C., a polyethylene terephthalate (PET) film can be used. The heat-resistant base material layer 11 may be made of a paper base material such as glassine paper, quality paper, or Japanese paper, or nonwoven fabric base material of cellulose, polyamide, polyester, aramid, or the like, or a metal film base material such as aluminum foil, SUS foil, or Ni foil. These materials may be stacked to form the base material layer 11 . The thickness of the base material layer 11 is at least 5 μm, preferably 10 to 200 μm, more preferably 25 to 100 μm, in order to prevent a rip and break. A thickness greater than or equal to 5 μm provides a good handling ability and a thickness less than or equal to 200 μm is advantageous in terms of cost. [Substrate Fixing Bond Layer 13 ] A bond used for the substrate fixing bond layer 13 may be the same resin as the resin used for the adhesive layer 12 or may be a material that has such adhesion strength that the substrate fixing bond layer 13 can be peeled away from the substrate or the base material layer 11 . Peeling of the heat-resistant adhesive sheet 2 away from the substrate 3 can be facilitated by heating if for example a blowing agent that is foamed by heat is added to the substrate fixing bond layer 13 . Instead of means that changes by heat, a component that forms cross-links under UV irradiation, for example, can be added to the substrate fixing bond layer 13 beforehand so that the substrate fixing bond layer 13 is cured, thereby reducing adhesion strength of the substrate fixing bond layer 13 . By such treatment, the adhesion strength of the substrate fixing bond layer 13 is reduced to separate the substrate 3 and the substrate fixing bond layer 13 from each other, or to separate the base material layer 11 and the substrate fixing bond layer 13 from each other, thereby removing the chips encapsulated with resin from the substrate 3 . [Flat, Smooth Peeling Sheet 10 ] The flat, smooth peeling sheet 10 is formed of a base material film having a peeling agent layer formed on one side of the base material film and is peeled to expose the adhesive layers on both sides before the heat-resistant adhesive sheet 2 for semiconductor device fabrication is used. The peeling agent layer contains a known peeling agent, such as a known fluorinated silicone resin-based peeling agent, a fluororesin peeling agent, a silicone resin-based peeling agent, polyvinyl alcohol-based resin, polypropylene-based resin, or long-chain alkyl compound, chosen according to the type of resin of the adhesive layer. The base material film may be any known film chosen from plastic films, such as polyether ketone, polyetherimide, polyarylate, polyethylene naphthalate, polyethylene, polypropylene, polybutene, polybutadiene, polymethylpentene, polyvinyl chloride, vinyl chloride copolymer, polyethylene terephthalate, polybutylene terephthalate, polyurethane, ethylene vinyl acetate copolymer, ionomeric resin, ethylene-(meth)acrylic acid copolymer, ethylene-(meth)acrylic acid ester copolymer, polystyrene, and polycarbonate films. [Fabrication of Heat-Resistant Adhesive Sheet 2 for Semiconductor Device Fabrication] According to the present invention, compositions prepared as described above can be used to form a heat-resistant adhesive sheet 2 for semiconductor device fabrication by any of method generally used for fabricating a multilayer structure. In one method, an adhesive in a solventless state or dissolved in a reactive solvent is applied on a base material film, then is dried by heating to form the heat-resistant adhesive sheet 2 for semiconductor device fabrication. Any of the heat-resistant base materials enumerated above is suitable as the material of the base material layer. An adhesive layer to be formed on the other side of the base material layer can be formed in the same way described above. Alternatively, an adhesive may be poured onto a peeling film or the like to form each individual film and the films may be stacked in sequence, or the application of an adhesive liquid described above may be combined with the stacking of individual films. Here, the solvent used is not limited particularly. Given that the material of the adhesive layer 12 has a good solubility, a ketone-based solvent such as methyl ethyl ketone is used preferably. Alternatively, the material may be an aquatic dispersion solution. The aquatic dispersion solution may be applied to the base material layer 11 and dried by heat. The process may be repeated to form the adhesive layer 12 , thereby forming the heat-resistant adhesive sheet 2 for semiconductor device fabrication. The heat-resistant adhesive sheet 2 for semiconductor device fabrication of the present invention includes the adhesive layer thus formed to a typical thickness of 1 to 50 μm on the base material layer. The heat-resistant adhesive sheet 2 for semiconductor device fabrication can be provided with an antistatic function as required. A method for providing an antistatic function to the heat-resistant adhesive sheet 2 for semiconductor device fabrication will be described. One method for providing the antistatic function is to add an antistatic agent or conductive filler to the adhesive layer 12 and the base material layer 11 . Another method is to apply an antistatic agent to the interface between the base material layer 11 and the adhesive layer 12 or the other surface of the base material layer 11 on which the adhesive layer 12 is not provided. The antistatic function can reduce static electricity build-up caused when the heat-resistant adhesive sheet 2 for semiconductor device fabrication is removed from the semiconductor device. The antistatic agent may be any agent that has the antistatic capability. Examples of the antistatic agent include surfactants such as acrylic-based ampholytic, acrylic-based cation, and maleic anhydride-styrene-based anion. Examples of the material for the antistatic layer include Bondeip PA, Bondeip PX, and Bondeip P (from Konishi Co., Ltd.). The conductive filler may be a conventional one, for example a metal such as Ni, Fe, Cr, Co, Al, Sb, Mo, Cu, Ag, Pt, or Au, or an alloy or oxide of any of these, or a carbon such as carbon black. These materials can be used either singly or in combination. The conductive filler may be powdery or fibrous filler. The heat-resistant adhesive sheet 2 for semiconductor device fabrication thus fabricated has an excellent heat resistance and a good demoldability from packages and therefore is suited for use in a semiconductor device manufacturing process. [Semiconductor Chip Bonding Step] The peeling sheets 10 are removed from both sides of the heat-resistant adhesive sheet 2 for semiconductor device fabrication, and the substrate fixing bond layer 13 is bonded onto the substrate so that the adhesive layer 12 is exposed at the top. Predetermined semiconductor chips 1 to be encapsulated with resin are placed and bonded on the adhesive layer 12 in a desired arrangement, thereby fixing the semiconductor chips 1 onto the adhesive layer 12 of the heat-resistant adhesive sheet 2 for semiconductor device fabrication. Here, the structure, shape, dimensions and the like of the semiconductor chips 1 are not limited to particular ones. [Encapsulation Step] Encapsulation resin 4 used in the encapsulation step in which the heat-resistant adhesive sheet 2 for semiconductor device fabrication of the present invention is used may be any known encapsulation resin such as epoxy resin. The melting temperature and curing temperature of resin powder or the curing temperature of a liquid resin are chosen by taking into account the heat resistance of the heat-resistant adhesive sheet 2 for semiconductor device fabrication. The heat-resistant adhesive sheet 2 for semiconductor device fabrication of the present invention is resistant to temperatures at which commonly-used encapsulation resins 4 cure and melt. The encapsulation step is performed in a mold by using any of the resins enumerated above for chip protection at a temperature from 170 to 180° C., for example. Then post-mold cure is performed before the heat-resistant adhesive sheet 2 for semiconductor device fabrication is peeled away. [Peeling Step] After the chips 1 fixed on the heat-resistant adhesive sheet 2 for semiconductor device fabrication on the substrate have been encapsulated with the resin, the structure is heated at a temperature in the range of 200 to 250° C. for 1 to 90 seconds (in the case of a hot plate) or 1 to 15 minutes (in the case of hot-air drier) to expand the substrate fixing bond layer 13 of the heat-resistant adhesive sheet 2 for semiconductor device fabrication, thereby reducing the bonding force between the substrate fixing bond layer 13 of the heat-resistant adhesive sheet 2 for semiconductor device fabrication and the substrate 3 to separate the heat-resistant adhesive sheet 2 and the substrate 3 from each other. Alternatively, the bonding force between the base material layer 11 of the heat-resistant adhesive sheet 2 for semiconductor device fabrication and the substrate fixing bond layer 13 is reduced to separate the base material layer 11 and the substrate fixing bond layer 13 from each other. Then the heat-resistant adhesive sheet 2 for semiconductor device fabrication is peeled away from the layer in which the chips are encapsulated with the resin. Alternatively, instead of separating the heat-resistant adhesive sheet 2 for semiconductor device fabrication and the substrate 3 , the chips 1 encapsulated with the resin 4 may be separated from the adhesive layer 12 of the heat-resistant adhesive sheet 2 for semiconductor device fabrication. [Electrode Forming Step] Then, screen printing or other method is used to form electrodes 5 in predetermined positions on the chips 1 on the side on which the heat-resistant adhesive sheet 2 was stacked and one surface of each of the chips 1 in the layer in which the chips 1 are encapsulated with the resin 4 is exposed. The electrodes 5 can be made of a known material. [Dicing Step] The layer in which the chips 1 are encapsulated with the resin 4 is fixed on a dicing sheet 8 on which preferably a dicing ring 7 is provided, then the layer is diced into packages with a dicing blade 6 for use in a typical dicing process. In doing so, if the chips 1 are not located in predetermined positions, the electrodes will be formed inaccurately, the locations of the chip 1 in each package will be inaccurate, or worse still, the dicing blade 6 can come into contact with the chips during the dicing. The heat-resistant adhesive sheet 2 for semiconductor device fabrication of the present invention can prevent displacements of the chips 1 during the encapsulation with the encapsulation resin 4 . Accordingly, the dicing step can be performed without the problems stated above and consequently packages in which the chips 1 are accurately positioned in the encapsulation resin 4 can be provided. The present invention will be described below with respect to working examples thereof. However, the present invention is not limited by these working examples. WORKING EXAMPLES [Measuring Method] Measurements and evaluations in working examples and comparative examples were made as follows. Initial adhesion strength to SUS: Peel adhesion strength to a SUS304BA plate at an angle of 180° at room temperature Adhesion strength to SUS304 after heating: Peel adhesion strength at an angle of 180° after attached to a SUS304BA plate and heated at 150° C. for 60 minutes Adhesion strength to SUS304 at 175° C.: Peel adhesion strength to a SUS304BA plate at 175° C. Peel strength from package: Peel adhesion strength at an angle of 180° when the adhesive sheet is peeled from the package Chip displacement: Displacement from the initial position of a chip measured with a digital microscope after package fabrication Adhesive deposit: The surface of the package was visually checked for adhesive deposits after the adhesive sheet was peeled off. The term “part” in the following description means “part by weight”. Working Example 1 3 pts.wt. of acrylic acid monomer as a component monomer was blended with 100 pts.wt. of butyl acrylate monomer to derive an acrylic-based copolymer. 3 pts.wt. of an epoxy-based cross-linking agent (Tetrad-C from Mitsubishi Gas Chemical Company, Inc.) and 5 pts.wt. of an isocyanate-based cross-linking agent (Coronate L from Nippon Polyurethane Industry Co., Ltd.) were blended with 100 pts.wt. of the acrylic-based copolymer to prepare an acrylic-based adhesive. 50 pts.wt. of a UV curable compound (UV-1700B from Nippon Synthetic Chemical Industry Co., Ltd.) and 3 pts.wt. of UV curing initiator (Irgacure 651 from NAGASE & CO., LTD.) were added to the acrylic-based adhesive to prepare an adhesive composition. Then, the adhesive composition was applied to a 25-μm-thick polyimide film (Kapton 100H from Du Pont-Toray Co., Ltd.) serving as the base material layer, and then dried to prepare a heat-resistant adhesive sheet for semiconductor device fabrication with an adhesive layer with a thickness approximately 10 μm. A 5 mm×5 mm Si wafer chip was placed on the heat-resistant adhesive sheet for semiconductor device fabrication, epoxy-based encapsulation resin powder (GE-7470LA from Nitto Denko Corporation) was sprinkled over the sheet and the wafer chip, and molded by heating at a temperature of 175° C. and a pressure of 3.0 kg/cm 2 for 2 minutes. Then, curing of the resin was accelerated by heating at 150° C. for 60 minutes (post-mold cure) to complete the package. Comparative Example 1 A package was fabricated in the same way as in Working Example 1, except that an adhesive sheet was peeled away before post-mold cure. Comparative Example 2 3 pts.wt. of acrylic acid monomer as a component monomer was blended with 100 pts.wt. of butyl acrylate monomer to prepare an acrylic-based copolymer. 0.6 pts.wt. of an epoxy-based cross-linking agent (Tetrad-C from Mitsubishi Gas Chemical Company, Inc.) and 2 pts.wt. of isocyanate-based cross-linking agent (Coronate L from Nippon Polyurethane Industry Co., Ltd.) were added to 100 pts.wt. of the acrylic-based copolymer to prepare an adhesive composition. The rest of the method for fabricating the package was the same as that in Working Example 1. Comparative Example 3 3 pts.wt. of acrylic acid monomer was added as a component monomer to 100 pts.wt. of butyl acrylate monomer to obtain an acrylic-based copolymer. 3 pts.wt. of an epoxy-based cross-linking agent (Tetrad-C from Mitsubishi Gas Chemical Company, Inc.) and 5 pts.wt. of an isocyanate-based cross-linking agent (Coronate L from Nippon polyurethane Industry Co., Ltd.) were added to 100 pts.wt. of the acrylic-based copolymer to prepare an adhesive composition. The rest of the method for fabricating the package was the same as that in Working Example 1. Results from Working Example 1 and Comparative Examples 1 to 3 are listed in Table 1 given below. TABLE 1 Compar- Compar- Compar- Working ative ative ative Unit example 1 example 1 example 2 example 3 Initial N/20 mm 0.71 0.71 0.80 0.24 adhesion strength to SUS304 Adhesion N/20 mm 0.35 0.35 2.62 1.22 strength to SUS304 after heating Peel N/20 mm 0.64 2.37 3.19 1.88 strength from package Chip mm 0.1 0.1 0.5 2.5 displacement Adhesive — Not Found Found Not found deposit found According to Working Example 1, a sufficient initial adhesion strength was achieved to suppress chip displacement. The adhesion strength was reduced because the adhesive was cured by heating during the post-mold cure. As a result, a light peel strength was achieved and a good package without an adhesive deposit after the peeling was successfully obtained. In Comparative Example 1, chip displacement was suppressed because of a sufficient initial adhesion strength as in Working Example 1. However, the peel strength was higher than that of Working Example 1 and adhesive deposits after the peeling occurred, because the sheet was peeled off with the adhesive being uncured before the post-mold cure. In Comparative Example 2, chip displacement was suppressed because the initial adhesion strength was high as in Working Example 1. However, since the adhesive used was not UV curable, adhesive strength reduction associated with curing of the adhesive by heating did not occur. Accordingly, the adhesion strength after heating was higher than the initial adhesion strength. Consequently, the peel strength from the package was high and adhesive deposit was left after the peeling. In Comparative Example 3, since the adhesive was not UV curable as in Comparative Example 2, the adhesion strength after heating was higher than the initial adhesion strength but was lower than that in Comparative Example 2 and therefore adhesive deposit did not occur. However, there was a large chip displacement because of the low initial adhesion strength. The results given above show that Working Example 1 was capable of providing a heat-resistant adhesive sheet for fabricating substrateless semiconductor packages that is capable of holding chips during resin encapsulation step while reducing adhesive deposit during peeling, because the adhesive layer was cured by heating after the resin encapsulation step. Working Example 2 42 parts of acrylonitrile-butadiene rubber (Nipol 1072) from Zeon Corporation), 53 parts of bisphenol A-type epoxy resin (Epikote 828 from Japan Epoxy Resin Co., Ltd., with a weight per epoxy equivalent of 190 g/eq), and 5 parts of imidazole (C1 1Z from Shikoku Chemicals Corporation) were blended and dissolved in an MEK solvent to a concentration of 35 wt % to prepare a bond solution. The bond solution was applied to a 35-μm-thick copper foil serving as a base material film, and was then dried at 150° C. for 3 minutes to form a bond layer having a bond thickness of 10 μm, thus forming a heat-resistant adhesive sheet for semiconductor device fabrication. A 3 mm×3 mm Si wafer chip was placed on the heat-resistant adhesive sheet for semiconductor device fabrication, epoxy-based encapsulation resin powder (GE-740LA from Nitto Denko Corporation) was sprinkled over the sheet and the wafer chip, and then molded by heating at a temperature of 175° C. under a pressure of 3.0 kg/cm 2 for 2 minutes. Then the structure was heated at 150° C. for 60 minutes to accelerate curing of the resin (post-mold cure) to complete a package. Working Example 3 24 parts of acrylonitrile-butadiene rubber (Nipol 1072) from Zeon Corporation), 65 parts of bisphenol A-type epoxy resin (Epikote 1002 from Japan Epoxy Resin Co., Ltd., with a weight per epoxy equivalent of 650 g/eq), 10 parts of phenol resin (P-180 from Arakawa Chemical Industries, Ltd.), and 1 part of triphenylphosphane (TPP from Hokko Chemical Industry) were blended and dissolved in an MEK solvent to a concentration of 35 wt % to prepare an adhesive solution. The adhesive solution was applied to a 35-μm-thick copper foil serving as the base material film, and was then dried at 150° C. for 3 minutes to form an adhesive layer having a thickness of 10 μm, thus forming a heat-resistant adhesive sheet for semiconductor device fabrication. The rest of the method for fabricating the package was the same as that in Working Example 2. Comparative Example 4 70 parts of acrylonitrile-butadiene rubber (Nipol 1072) from Zeon Corporation), 28 parts of bisphenol A-type epoxy resin (Epikote 828 from Japan Epoxy Resin Co., Ltd., with a weight per epoxy equivalent of 190 g/eq), and 2 parts of imidazole (C1 1Z from Shikoku Chemicals Corporation) were blended and were dissolved in an MEK solvent to a concentration of 35 wt % to prepare an adhesive solution. The adhesive solution was applied to a 35-μm-thick copper foil serving as a base material film and was then dried at 150° C. for 3 minutes to form an adhesive layer having an adhesive thickness of 10 μm, thus forming a heat-resistant adhesive sheet for semiconductor device fabrication. The rest of the method for fabricating the package was the same as that in Working Example 2. Comparative Example 5 10 parts of acrylonitrile-butadiene rubber (Nipol 1072) from Zeon Corporation), 79 parts of bisphenol A-type epoxy resin (Epikote 1002 from Japan Epoxy Resin Co., Ltd., with a weight per epoxy equivalent of 650 g/eq), 10 parts of phenol resin (P-180 from Arakawa Chemical Industries, Ltd.), and 1 part of triphenylphosphane (TPP from Hokko Chemical Industry) were blended and dissolved in an MEK solvent to a concentration of 35 wt % to prepare an adhesive solution. The adhesive solution was applied to a 35-μm-thick copper foil serving as the base material film, and was then dried at 150° C. for 3 minutes to form an adhesive layer having an adhesive thickness of 10 μm, thus forming a heat-resistant adhesive sheet for semiconductor device fabrication. The rest of the method for fabricating a package was the same as that in Working Example 2. Results from Working Examples 2 and 3 and Comparative Examples 4 and 5 are listed in Table 2 given below. TABLE 2 Compar- Compar- Working Working ative ative Unit example 2 example 3 example 4 example 5 Initial N/20 mm 1.22 1.14 3.12 0.27 adhesion strength to SUS304 Peel N/20 mm 0.82 0.72 2.32 0.21 strength from package Chip mm 0.1 0.1 1.0 2.9 displacement Adhesive — Not Not Found Not found deposit found found As apparent from Table 2, the adhesive sheets of Working Example 2 and 3 of the present invention had an excellent demoldability and did not leave adhesive deposits. Furthermore, the adhesive sheets had sufficient initial adhesion strength to suppress chip displacement. In contrast, the adhesive sheet of Comparative Example 4 with a large amount of rubber component had sufficient initial adhesion strength but caused chip displacement during resin encapsulation because the adhesive layer was soft. In addition, the adhesive sheet had poor elasticity after cured and left adhesive deposit. The adhesive sheet of Comparative Example 5 with a small amount of rubber component had insufficient initial adhesion strength and therefore caused chip displacement during resin encapsulation. Working Example 4 1.0 parts of epoxy-based cross-linking agent (Tetrad-C from Mitsubishi Gas Chemical Company, Inc.), 5 parts of a rosin phenolic tackifier, 50 parts of thermally expandable microspheres that form at 200° C. and toluene were uniformly blended and dissolved with 100 parts of a copolymer including ethyl acrylate-butyl acrylate-acrylic acid (20 parts-80 parts-10 parts) to prepare a coating solution. Then, an addition reactive silicone adhesive (SD-4587L from Dow Corning Toray Co., Ltd.) was applied to a 25-μm-thick polyimide film (Kapton 100H from Du Pont-Toray Co., Ltd.) to a thickness of 5 μm and dried. The film was used as a base material layer and the adhesive composition described above was applied to the surface of the film on which the silicone adhesive was not applied and was then dried to form a heat-resistant adhesive sheet for semiconductor device fabrication with the adhesive layer having a thickness of approximately 40 μm. The thermally expandable adhesive surface of the heat-resistant adhesive sheet for semiconductor device fabrication was fixed on a flat, smooth platform by pressure and a 5 mm×5 mm Si wafer chip was placed on the silicone adhesive surface. Epoxy-based resin powder (GE-7470LA from Nitto Denko Corporation) was sprinkled over the sheet and the wafer chip, and was then molded by heating at a temperature of 175° C. under a pressure of 400 kpa for 2 minutes. Then the structure was heated at 150° C. for 60 minutes to accelerate the curing of the resin (post-mold cure) to form a package. Working Example 5 A package was fabricated in the same way as in Working Example 4, except that 1.0 parts of epoxy-based cross-linking agent (Tetrad-C from Mitsubishi Gas Chemical Company, Inc.), 5 parts of a rosin phenolic tackifier, and toluene were uniformly blended and dissolved with 100 parts of copolymer including ethyl acrylate-butyl acrylate-acrylic acid (20 parts-80 parts-10 parts) to prepare a coating solution. Comparative Example 6 A package was fabricated in the same way as in Working Example 4, except that instead of the silicone adhesive of the heat-resistant adhesive sheet for semiconductor device fabrication, an adhesive composition was applied to a polyimide film (Kapton 100H from Du Pont-Toray Co., Ltd.). The adhesive composition was prepared by adding 0.6 pts.wt. of an epoxy-based cross-linking agent (Tetrad-C from Mitsubishi Gas Chemical Company, Inc.) and 2 pts.wt. of isocyanate-based cross-linking agent (Coronate L from Nippon Polyurethane Industry Co., Ltd.) to 100 pts.wt. of acrylic-based copolymer including 3 pts.wt. of acrylic acid monomer as a component monomer for 100 pts.wt. of butyl acrylate monomer. Results from Working Examples 4 and 5 and Comparative Example 6 are listed in Table 3 given below. TABLE 3 Working Working Comparative Unit Example 4 Example 5 Example 6 Adhesion N/20 mm 0.67 0.72 0.15 strength to SUS304 at 175° C. Peel N/20 mm 0.34 0.48 2.46 strength from package Chip mm 0.1 0.1 2.5 displacement Adhesive — Not found Not found Found deposit In Working Example 4, the adhesive sheet had a sufficient adhesive strength to suppress chip displacement during mold. In addition, light package peel strength was successfully achieved by using the silicone adhesive's characteristic of low adhesive strength to encapsulation resin, and therefore a good package without adhesive deposit was obtained. In Working Example 5, the adhesive sheet had high adhesion strength at 175° C. to suppress chip displacement. In addition, light package peel strength was successfully achieved and therefore a good package without adhesive deposit was obtained. In Comparative Example 6, chip displacement was not suppressed because of reduction in adhesive strength in a high temperature range, which is a characteristic of acrylic adhesives. The result shows that chip displacement cannot be suppressed unless the adhesive layer is a silicone adhesive layer, even though the base material layer has a sufficiently low thermal expansion coefficient. In addition, peel strength was higher than that in the working examples because of high adhesion strength to the encapsulation resin, which is polyfunctional. Consequently, adhesive deposits were left after peeling.

The present invention is intended to solve the following problems with a method for fabricating a substrateless semiconductor package using an adhesive sheet as a temporary fixing supporter. A chip can be displaced from a specified position by pressure during resin encapsulation because the chip is not properly held by the adhesive sheet. If such displacement occurs, the relative positional relationship between the chip and an interconnect to be connected to a specified position in a subsequent wiring step also changes by the displacement of the chip from the specified position. Another problem is that if adhesive deposits occur during peeling of the adhesive sheet and the surface of a package is contaminated with the adhesive deposits, adhesive components left on the surface of the chip can inhibit connection between the interconnect and the chip in a subsequent wiring step. To solve these problems, the present invention provides an adhesive sheet for semiconductor device fabrication that is attached to a substrateless semiconductor chip when the chip is encapsulated with resin. The adhesive sheet includes a base material layer and an adhesive layer. The adhesive layer has a specific adhesion strength and peel strength.

CROSS-REFERENCE TO RELATED APPLICATIONS [0001] This application is a continuation of U.S. patent application Ser. No. 13/339,236, filed Dec. 28, 2011, which claims priority to U.S. Provisional Patent Application No. 61/428,114, filed Dec. 29, 2010. U.S. patent application Ser. No. 13/339,236 is a continuation-in-part of U.S. patent application Ser. No. 11/888,009, filed Jul. 31, 2007, now U.S. Pat. No. 8,252,036, which claims priority to U.S. Provisional Patent Application Nos. 60/834,401 filed Jul. 31, 2006 and 60/834,627 filed Aug. 1, 2006. U.S. patent application Ser. No. 13/339,236 is also a continuation-in-part of U.S. patent application Ser. No. 12/822,291, filed Jun. 24, 2010, now U.S. Pat. No. 9,408,607, which claims priority to U.S. Provisional Patent Application No. 61/222,646, filed Jul. 2, 2009. All of the above-referenced patents and applications are incorporated by reference herein in their entireties. STATEMENT REGARDING FEDERALLY SPONSORED RESEARCH OR DEVELOPMENT [0002] Not Applicable FIELD OF THE INVENTION [0003] The present invention relates to the field of surgical implant devices and methods for their manufacture and use. Among the exemplary embodiments of the present invention are improvements in sealing and retention medical devices particularly applicable to vascular surgery and the treatment of aneurysms or other luminal defects in other anatomic conduits, such as sealing and retention of replacement heart valves. BACKGROUND OF THE INVENTION [0004] Medical and surgical implants are placed often in anatomic spaces where it is desirable for the implant to conform to the unique anatomy of the targeted anatomic space and secure a seal therein, preferably without disturbing or distorting the unique anatomy of that targeted anatomic space. [0005] While the lumens of most hollow anatomic spaces are ideally circular, in fact, the cross-sectional configurations of most anatomic spaces are, at best, ovoid, and may be highly irregular. Such lumenal irregularity may be due to anatomic variations and/or to pathologic conditions that may change the shape and topography of the lumen and its associated anatomic wall. Examples of anatomic spaces where such implants may be deployed include, but are not limited to, blood vessels, the heart, other vascular structures, vascular defects (such as thoracic and abdominal aortic aneurysms), the trachea, the oropharynx, the esophagus, the stomach, the duodenum, the ileum, the jejunum, the colon, the rectum, ureters, urethras, fallopian tubes, biliary ducts, pancreatic ducts, or other anatomic structures containing a lumen used for the transport of gases, blood, or other liquids or liquid suspensions within a mammalian body. [0006] For a patient to be a candidate for existing endograft methods and technologies, to permit an adequate seal, a proximal neck of, ideally, at least 12 mm of normal aorta must exist downstream of the left subclavian artery for thoracic aortic aneurysms or between the origin of the most inferior renal artery and the origin of the aneurysm in the case of abdominal aneurysms. Similarly, ideally, at least 12 mm of normal vessel must exist distal to the distal extent of the aneurysm for an adequate seal to be achieved. [0007] Migration of existing endografts has also been a significant clinical problem, potentially causing leakage and profusion of aneurysms and/or compromising necessary vascular supplies to arteries such as the carotid, subclavian, renal, or internal iliac vessels. This problem only has been addressed partially by some existing endograft designs, in which barbs or hooks have been incorporated to help retain the endograft at its intended site. However, most existing endograft designs are solely dependent on radial force applied by varying length of stent material to secure a seal against the recipient vessel walls. [0008] Because of the limitations imposed by existing vascular endograft devices and endovascular techniques, a significant number of abdominal and thoracic aneurysms repaired in the U.S. are still managed though open vascular surgery, instead of the lower morbidity of the endovascular approach. [0009] Pre-sizing is required currently in all prior art endografts. Such pre-sizing based on CAT-scan measurements is a significant problem. This leads, many times, to mis-sized grafts. In such situations, more grafts segments are required to be placed, can require emergency open surgery, and can lead to an unstable seal and/or migration. Currently there exists no endograft that can be fully repositioned after deployment. [0010] Thus, a need exists to overcome the problems with the prior art systems, designs, and processes as discussed above. SUMMARY OF THE INVENTION [0011] The invention provides surgical implant devices and methods for their manufacture and use that overcome the hereinafore-mentioned disadvantages of the heretofore-known devices and methods of this general type and that provide such features with improvements that increase the ability of such an implant to be precisely positioned and sealed, with better in situ accommodation to the local anatomy of the targeted anatomic site. The invention provide an adjustment tool that can remotely actuate an adjustment member(s) that causes a configuration change of a portion(s) of an implant, which configuration change includes but is not limited to diameter, perimeter, shape, and/or geometry or a combination of these, to create a seal and provide retention of an implant to a specific area of a target vessel or structure. [0012] One exemplary aspect of the present invention is directed towards novel designs for endovascular implant grafts, and methods for their use for the treatment of aortic aneurysms and other structural vascular defects. An endograft system for placement in an anatomic structure or blood vessel is disclosed in which an endograft implant comprises, for example, a non-elastic tubular implant body with at least an accommodating proximal end. Accommodating, as used herein, is the ability to vary a configuration in one or more ways, which can include elasticity, expansion, contraction, and changes in geometry. Both or either of the proximal and distal ends in an implant according to the present invention further comprise one or more circumferential expandable sealable collars and one or more expandable sealing devices, capable of being expanded upon deployment to achieve the desired seal between the collar and the vessel's inner wall. Exemplary embodiments of such devices can be found in co-pending U.S. patent application Ser. No. 11/888,009, filed Jul. 31, 2007, and Ser. No. 12/822,291, filed Jun. 24, 2010, which applications have been incorporated herein in their entireties. Further embodiments of endovascular implants according to the present invention may be provided with retractable retention tines or other retention devices allowing an implant to be repositioned before final deployment. In other embodiments, the implant can be repositioned after final deployment. An endograft system according to the present invention further comprises a delivery catheter with an operable tubular sheath capable of housing a folded or compressed endograft implant prior to deployment and capable of retracting or otherwise opening in at least its proximal end to allow implant deployment. The sheath is sized and configured to allow its placement via a peripheral arteriotomy site, and is of appropriate length to allow its advancement into the aortic valve annulus, ascending aorta, aortic arch, and thoracic or abdominal aorta, as required for a specific application. [0013] While some post-implantation remodeling of the aortic neck proximal to an endovascular graft (endograft) has been reported, existing endograft technology does not allow for the management of this condition without placement of an additional endograft sleeve to cover the remodeled segment. [0014] Exemplary endografts of the present invention as described herein allow for better accommodation by the implant of the local anatomy, using a self-expandable or compressible gasket for the sealing interface between the endograft collar and the recipient vessel's inner wall. Furthermore, exemplary endografts of the present invention as disclosed herein are provided with a controllably releasable disconnect mechanism that allows remote removal of an adjustment tool and locking of the retained sealable mechanism after satisfactory positioning and sealing of the endograft. In some exemplary embodiments according to the present invention, the controllably releasable disconnect mechanism may be provided in a manner that allows post-implantation re-docking of an adjustment member to permit post-implantation repositioning and/or resealing of an endograft subsequent to its initial deployment. [0015] In other exemplary applications encompassed by the present invention, improved devices for sealing other medical devices such as vascular cannulae may be provided. The present invention further includes novel designs for vascular cannulae to be used when bi-caval cannulation of the heart is indicated, eliminating the need to perform circumferential caval dissection and further reducing the tissue trauma caused by prior art balloon or other bypass cannulae. While the vascular cannulae of the present invention are inserted and positioned by a surgeon in the standard fashion, the need for circumferential dissection of the cavae and tourniquet placement is obviated. After the vascular cannulae of the present invention are positioned and secured with purse string sutures, the surgeon deploys the adjustable sealing devices of the cannulae by turning an adjustment tool or torque wire. Once the sealing devices are deployed, all of the venous return is diverted. The sealing devices deploy around the distal ends of the cannulae and allow blood to flow through the lumen of the cannulae, but not around the sealing devices. Use of these cannulae minimizes the chance of caval injury by eliminating the need for circumferential dissection. Additionally, the configuration of the adjustable sealing device in relation to the cannula is such that the adjustable sealing device is “flush” with the cannula so that no acute change in diameter exists along the external surface of the cannula, which serves to avoid tissue trauma during insertion and withdrawal into and out of bodily structures. [0016] The present invention addresses several major problems presented by existing designs for balloon cannulae. In various exemplary embodiments according to the present invention, the lumens are configured such that a cannula with an adjustable sealing device can be deployed without compromising either the flow within the principle lumen of the cannula or the seal between the cannula and the structure within which the cannula lies. Moreover, a disclosed example of a cannula according to the present invention is provided with a trough within the cannula body at its distal end in which the adjustable sealing device member lies such that, when undeployed during insertion and withdrawal, there is a smooth interface between the external cannula wall and the undeployed sealing device, allowing for smoother, easier, and safer insertion and withdrawal. [0017] Moreover, existing designs for balloon cannulae are unable to provide a truly symmetrical placement of an inflated balloon around a central lumen of standard diameter. The asymmetry that results with conventional balloon inflation is sufficient to displace the lumen from the true center of the endovascular lumen in which the balloon cannula is placed, resulting in unpredictable and suboptimal flow characteristics therethrough. The altered hemodynamics of such flow with an existing balloon cannula increases the likelihood of intimal vascular injury and clot or plaque embolization. Vascular cannulae of the present invention achieve the surprising result of having the flow characteristics of a non-balloon cannula by maintaining the preferred laminar flow characteristics of a circular main lumen of consistent diameter, positioned and maintained in or near the center of vascular flow by an adjustable sealing device originally provided within a recessed trough in the exterior wall of the cannula, with accessory lumens contained within an externally circular cannular wall. This allows for better seal, less vascular trauma, and easier vascular ingress and egress. [0018] In addition, vascular cannulae according to the present invention may be provided with retractable stabilizing elements to anchor the inflated balloon within a vessel lumen during use. Such stabilizing elements further make use of the trough within the cannula body, with the stabilizing elements retracting into this trough during insertion and removal, allowing for smooth and trauma-free entry and egress of the cannula. [0019] Certain aspects of the present invention are directed towards novel designs for sealable endovascular implant grafts, and methods for their use for the treatment of aortic aneurysms and other structural vascular defects or for heart valve replacements. Various embodiments as contemplated within the present invention may include any combination of exemplary elements as disclosed herein or in the co-pending patent applications referenced above. [0020] In an exemplary embodiment according to the present invention, a sealable vascular endograft system for placement in a vascular defect is provided, comprising an elongated main implant delivery catheter with an external end and an internal end for placement in a blood vessel with internal walls. In such an exemplary embodiment, the main implant delivery catheter further comprises a main implant delivery catheter sheath that may be openable or removable at the internal end and a main implant delivery catheter lumen containing within a compressed or folded endovascular implant. Further, in such an exemplary embodiment, an endovascular implant comprises a non-elastic tubular implant body with an accommodating proximal end terminating in a proximal sealable circumferential collar that may be expanded by the operator to achieve a fluid-tight seal between the proximal sealable circumferential collar and the internal walls of the blood vessel proximal to the vascular defect. Moreover, in such an exemplary embodiment, an endovascular implant may further comprises a non-elastic tubular implant body with an accommodating distal end terminating in a distal sealable circumferential collar controlled by a distal variable sealing device, which may be expanded by the operator to achieve a fluid-tight seal between the distal sealable circumferential collar and the internal walls of the blood vessel distal to the vascular defect. [0021] In a further exemplary embodiment according to the present invention, an implant interface is provided for a sealable attachment of an implant to a wall within the lumen of a blood vessel or other anatomic conduit. [0022] In a yet further exemplary embodiment according to the present invention, an implant gasket interface is provided for a sealable attachment of an implant to a wall within the lumen of a blood vessel or other anatomic conduit, wherein the sealable attachment provides for auto-adjustment of the seal while maintaining wall attachment to accommodate post-implantation wall remodeling. [0023] Still other exemplary embodiments of endografts and endograft delivery systems according to the present invention serve as universal endograft cuffs, being first placed to offer their advantageous anatomic accommodation capabilities, and then serving as a recipient vessel for other endografts, including conventional endografts. [0024] Furthermore, exemplary embodiments of endografts and endograft delivery systems according to the present invention may be provided with a mechanism to permit transfer of torque or other energy from a remote operator to an adjustment member comprising a sealable, adjustable circumferential assembly controlled by an adjustment tool, which may be detachable therefrom and may further cause the assembly to lock upon detachment of the tool. In some exemplary embodiments of the present invention, the variable sealing device may be provided with a re-docking element that may be recaptured by subsequent operator interaction, allowing redocking and repositioning and/or resealing of the endograft at a time after its initial deployment. [0025] Moreover, the various exemplary embodiments of the present invention as disclosed herein may constitute complete endograft systems, or they may be used as components of a universal endograft system as disclosed in co-pending patent applications that may allow the benefits of the present invention to be combined with the ability to receive other endografts. [0026] Finally, the present invention encompasses sealable devices that may be used in other medical devices such as adjustable vascular cannulas or other medical or surgical devices or implants, such as aortic valves. [0027] With the foregoing and other objects in view, there is provided, in accordance with the invention, a surgical implant including an implant body and a selectively adjustable assembly attached to the implant body, having adjustable elements, and operable to cause a configuration change in a portion of the implant body and, thereby, permit implantation of the implant body within an anatomic orifice to effect a seal therein under normal physiological conditions. [0028] The preceding description is presented only as an exemplary application of the devices and methods according to the present invention. [0029] Although the invention is illustrated and described herein as embodied in surgical implant devices and methods for their manufacture and use, it is, nevertheless, not intended to be limited to the details shown because various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims. Additionally, well-known elements of exemplary embodiments of the invention will not be described in detail or will be omitted so as not to obscure the relevant details of the invention. [0030] Additional advantages and other features characteristic of the present invention will be set forth in the detailed description that follows and may be apparent from the detailed description or may be learned by practice of exemplary embodiments of the invention. Still other advantages of the invention may be realized by any of the instrumentalities, methods, or combinations particularly pointed out in the claims. [0031] Other features that are considered as characteristic for the invention are set forth in the appended claims. As required, detailed embodiments of the present invention are disclosed herein; however, it is to be understood that the disclosed embodiments are merely exemplary of the invention, which can be embodied in various forms. Therefore, specific structural and functional details disclosed herein are not to be interpreted as limiting, but merely as a basis for the claims and as a representative basis for teaching one of ordinary skill in the art to variously employ the present invention in virtually any appropriately detailed structure. Further, the terms and phrases used herein are not intended to be limiting; but rather, to provide an understandable description of the invention. While the specification concludes with claims defining the features of the invention that are regarded as novel, it is believed that the invention will be better understood from a consideration of the following description in conjunction with the drawing figures, in which like reference numerals are carried forward. BRIEF DESCRIPTION OF THE DRAWINGS [0032] The accompanying figures, where like reference numerals refer to identical or functionally similar elements throughout the separate views, which are not true to scale, and which, together with the detailed description below, are incorporated in and form part of the specification, serve to illustrate further various embodiments and to explain various principles and advantages all in accordance with the present invention. Advantages of embodiments of the present invention will be apparent from the following detailed description of the exemplary embodiments thereof, which description should be considered in conjunction with the accompanying drawings in which: [0033] FIG. 1 is a fragmentary, perspective view of an exemplary embodiment of a proximal aspect of a selectively expandable and contractable endograft according to the present invention with the endograft in a relatively expanded form; [0034] FIG. 2 is a fragmentary, perspective view of the selectively expandable and contractable endograft of FIG. 1 with the endograft in a relatively contracted form; [0035] FIG. 3 is a fragmentary, perspective view of another exemplary embodiment of a proximal aspect of an endograft according to the present invention further incorporating a lattice structure; [0036] FIG. 4A is a fragmentary, perspective view of the endograft of FIG. 3 with the endograft in a relatively contracted form; [0037] FIG. 4B is a fragmentary, perspective view of the endograft of FIG. 3 with the endograft in a partially expanded form; [0038] FIG. 4C is a fragmentary, perspective view of the endograft of FIG. 3 with the endograft in a fully expanded form; [0039] FIG. 5A is a fragmentary, partially hidden, perspective view of an exemplary embodiment of a microcylinder locking mechanism with an associated adjustment tool prior to engagement of the microcylinder locking mechanism by the adjustment tool; [0040] FIG. 5B is a fragmentary, partially hidden, perspective view of the microcylinder locking mechanism and adjustment tool of FIG. 5B with engagement of the microcylinder locking mechanism by the adjustment tool; [0041] FIG. 5C is a fragmentary, partially hidden, perspective view of an exemplary embodiment of the microcylinder locking mechanism and adjustment tool of FIG. 5B after adjustment and disengagement of the adjustment tool from the microcylinder locking mechanism; [0042] FIG. 6A is an axial cross-sectional view of the microcylinder and guide bullet along section line A-A of FIG. 5A with tines captures in striations of the microcylinder; [0043] FIG. 6B is an axial cross-sectional view of the adjustment tool along section line B-B of FIG. 5A ; [0044] FIG. 6C is an axial cross-sectional view of the microcylinder along section line C-C of FIG. 5B ; [0045] FIG. 6D is an axial cross-sectional view of the microcylinder, the guide bullet, and the tool sheath along section line D-D of FIG. 5B without the adjustment member with the tines removed from the microcylinder by the adjustment tool; [0046] FIG. 6E is an axial cross-sectional view of another exemplary embodiment of a microcylinder locking mechanism and adjustment tool sheath according to the invention where the adjustment tool also has striations having a rectangular cross-sectional shape and has a smooth exterior; [0047] FIG. 6F is an axial cross-sectional view of yet another exemplary embodiment of a microcylinder locking mechanism according to the invention in which the microcylinder has striations with a triangular cross-sectional shape and with the tines caught in the striations of the microcylinder; [0048] FIG. 6G is an axial cross-sectional view of the microcylinder locking mechanism of FIG. 6F and an adjustment tool according to the invention in which the tines are removed from the microcylinder by the adjustment tool; [0049] FIG. 7A is a longitudinal, partial cross-sectional view of an exemplary embodiment of an adjustment control locking mechanism according to the present invention with a controllable catch mechanism disengaged; [0050] FIG. 7B is a longitudinal, partial cross-sectional view of the adjustment control locking mechanism of FIG. 7A with the controllable catch mechanism engaged. [0051] FIG. 8A is a fragmentary, partially hidden, perspective view of an exemplary embodiment of a microcylinder locking mechanism according to the invention with internal locking tines of unequal length and with an associated adjustment tool sheath prior to engagement of the microcylinder locking mechanism by the adjustment tool sheath; [0052] FIG. 8B is a fragmentary, partially hidden, perspective view of the microcylinder locking mechanism and adjustment tool sheath of FIG. 7A with engagement of the microcylinder locking mechanism by the adjustment tool sheath; [0053] FIG. 8C is a fragmentary, partially hidden, perspective view of the microcylinder locking mechanism and adjustment tool sheath of FIG. 7B after adjustment and disengagement of the microcylinder locking mechanism with the adjustment tool sheath. [0054] FIG. 9A is an axial cross-sectional view of retention tines sheathed by an expanded compressible foam gasket in an exemplary endograft according to the present invention with the tines in a non-extended state; [0055] FIG. 9B is a fragmentary, perspective view of the retention tines of FIG. 9A exposed and deployed through a compressible foam gasket by an expanded sealable collar in an exemplary endograft according to the present invention; [0056] FIG. 10A is a fragmentary, axial cross-sectional view of an exemplary endovascular interface cuff according to the present invention, in which the interface cuff has been positioned over an endovascular guidewire to a desired recipient site in the aorta proximal to an aortic aneurysm sac but has not been expanded therein; [0057] FIG. 10B is a fragmentary, transverse cross-sectional view of the interface cuff of FIG. 10A ; [0058] FIG. 11A is a fragmentary, axial cross-sectional view of the interface cuff of FIG. 10A , with expansion of the endovascular interface cuff in the aorta to achieve a seal and with retention tine engagement of the aortic wall in the desired recipient site proximal to the aortic aneurysm sac at the level of A-A′; [0059] FIG. 11B is a fragmentary, transverse cross-sectional view of the interface cuff of FIG. 11A ; [0060] FIG. 12 is a fragmentary, axial cross-sectional view of the interface cuff of FIG. 10A with delivery of an endograft secured within the rigid cuff of the interface cuff; [0061] FIG. 13 is a fragmentary, axial cross-sectional view of the interface cuff of FIG. 12 with the guidewire removed and with the adjustment tool detached and removed; [0062] FIG. 14A is a fragmentary, perspective view of an exemplary embodiment of an actively controllable endograft according to the present invention in which a latticework external to the lumen of an endograft can be radially displaced by controlled rotation of an adjustment member, the lattice structure being in a contracted state; [0063] FIG. 14B is a fragmentary, perspective view of the actively controllable endograft of FIG. 14A in which the lattice structure is in an expanded state; [0064] FIG. 15A is a side perspective view of an exemplary embodiment of an adjustable vascular cannula according to the present invention; [0065] FIG. 15B is a side perspective and partially hidden view of the adjustable vascular cannula of FIG. 15A within a recipient blood vessel with an adjustable seal device in a non-deployed, contracted position; and [0066] FIG. 15C is a side perspective and partially hidden view of the adjustable vascular cannula of FIG. 15B with the adjustable seal device in a deployed, expanded position. DETAILED DESCRIPTION OF THE INVENTION [0067] As required, detailed embodiments of the present invention are disclosed herein; however, it is to be understood that the disclosed embodiments are merely exemplary of the invention, which can be embodied in various forms. Therefore, specific structural and functional details disclosed herein are not to be interpreted as limiting, but merely as a basis for the claims and as a representative basis for teaching one skilled in the art to variously employ the present invention in virtually any appropriately detailed structure. Further, the terms and phrases used herein are not intended to be limiting; but rather, to provide an understandable description of the invention. While the specification concludes with claims defining the features of the invention that are regarded as novel, it is believed that the invention will be better understood from a consideration of the following description in conjunction with the drawing figures, in which like reference numerals are carried forward. [0068] Alternate embodiments may be devised without departing from the spirit or the scope of the invention. Additionally, well-known elements of exemplary embodiments of the invention will not be described in detail or will be omitted so as not to obscure the relevant details of the invention. [0069] Before the present invention is disclosed and described, it is to be understood that the terminology used herein is for the purpose of describing particular embodiments only and is not intended to be limiting. The terms “a” or “an”, as used herein, are defined as one or more than one. The term “plurality,” as used herein, is defined as two or more than two. The term “another,” as used herein, is defined as at least a second or more. The terms “including” and/or “having,” as used herein, are defined as comprising (i.e., open language). The term “coupled,” as used herein, is defined as connected, although not necessarily directly, and not necessarily mechanically. [0070] Relational terms such as first and second, top and bottom, and the like may be used solely to distinguish one entity or action from another entity or action without necessarily requiring or implying any actual such relationship or order between such entities or actions. The terms “comprises,” “comprising,” or any other variation thereof are intended to cover a non-exclusive inclusion, such that a process, method, article, or apparatus that comprises a list of elements does not include only those elements but may include other elements not expressly listed or inherent to such process, method, article, or apparatus. An element proceeded by “comprises . . . a” does not, without more constraints, preclude the existence of additional identical elements in the process, method, article, or apparatus that comprises the element. [0071] As used herein, the term “about” or “approximately” applies to all numeric values, whether or not explicitly indicated. These terms generally refer to a range of numbers that one of skill in the art would consider equivalent to the recited values (i.e., having the same function or result). In many instances these terms may include numbers that are rounded to the nearest significant figure. [0072] Herein various embodiments of the present invention are described. In many of the different embodiments, features are similar. Therefore, to avoid redundancy, repetitive description of these similar features may not be made in some circumstances. It shall be understood, however, that description of a first-appearing feature applies to the later described similar feature and each respective description, therefore, is to be incorporated therein without such repetition. [0073] Described now are exemplary embodiments of the present invention. Referring now to the figures of the drawings in detail and, first, particularly to FIG. 1 thereof, there is shown a perspective view of an exemplary embodiment of the proximal aspect of a sealable endograft system 1000 according to the present invention, in which the endograft is in a relatively expanded form. FIG. 2 is a perspective view of the embodiment of the proximal aspect of a sealable endograft system 1000 according to the present invention of FIG. 1 , showing the endograft in a relatively contracted form. This exemplary endograft system 1000 has the ability to be selectively expanded and contracted to a diameter selected by the implanting physician. In general, the endograft system 1000 has, along its intermediate extent and, possibly, also at its distal portion (at the downstream end of the prosthesis), a relatively constant diameter portion. At its proximal portion (at the upstream end of the prosthesis), the endograft system 1000 is able to impart a configuration change to selectively adjustable portion of the implant. Features of the inventive controllable endograft system 1000 are described in further detail in U.S. patent application Ser. No. 11/888,009, filed Jul. 31, 2007, and Ser. No. 12/822,291, filed Jun. 24, 2010, which have been incorporated herein and detail of which is not replicated herein for the sake of brevity. [0074] The exemplary sealable endograft system 1000 shown in FIGS. 1 and 2 comprises a hollow tubular endograft body 1005 having an accommodating proximal cuff 1010 and an intermediate, substantially rigid, tubular member 1015 . The distal end of such an endograft (not shown in FIGS. 1 and 2 ) may be any or all of accommodating, elastic, rigid, stent-laden, or even replicate the proximal end, depending upon the various exemplary embodiments according to the present invention. A selectively adjustable circumferential assembly 1020 is disposed at the proximal cuff 1010 . Contained in one exemplary embodiment of the circumferential assembly 1020 is a circumferential channel enclosing an adjustment member 1025 (indicated only diagrammatically with a solid line). The adjustment member 1025 causes the expansion/contraction of the accommodating proximal cuff 1010 by looping around the perimeter and by being lengthened or shortened, respectively. The adjustment member 1025 , for example, interacts with a control device 1030 that is operable to cause an increase or decrease in the circumference of the circumferential loop 1025 by the application of rotational torque to the distal aspect of an adjustment tool 1035 emerging from the control device 1030 . The adjustment member 1025 can be integral with the adjustment tool 1035 in an exemplary embodiment of the circumferential assembly 1020 , or can be removable as shown, for example, in FIG. 10A . [0075] Such an adjustment member 1025 may take many forms in the present invention. In one exemplary embodiment according to the present invention, the adjustment member 1025 is a micro-threaded cable that is fixed at one end to the control device 1030 , which is in the form of a microcylinder, and the adjustment tool 1035 threads through a threaded aspect of the microcylinder 1030 in order to effect a change in the circumference of the proximal cuff 1010 . A forwardly imposed torque on the adjustment tool 1035 cause expansion of the adjustment tool 1035 . Expansion of the adjustment member 1025 in its circumferential extent has the effect of expanding the proximal aspect of the sealable endograft system 1000 to allow for precise sealing of the sealable endograft system 1000 within a recipient blood vessel such as the aorta (not shown in FIG. 1 or 2 ). Conversely, reverse torque on the adjustment tool 1035 has the effect of decreasing the circumference of the circumferential loop of the adjustment member 1025 and, thus, contracting the proximal aspect of the sealable endograft system 1000 , allowing for re-positioning as needed. In FIGS. 1 and 2 , the adjustment tool 1035 may extend distally through the lumen of the sealable endograft system 1000 . Alternatively, the adjustment tool 1035 may extend distally through a separate lumen provided in the sealable endograft system 1000 (not shown in FIG. 1 or 2 ). [0076] FIGS. 3 and 4A to 4C are perspective views of yet another exemplary embodiment of a proximal aspect of a sealable endograft system 1000 according to the present invention that further incorporates a stent or lattice structure 1041 (which, in another embodiment, can be a compressible foam gasket). The lattice structure 1041 is provided with a lattice interruption 1045 to allow for variations in the circumference of the proximal aspect of the endograft. This lattice interruption 1045 may take the form of a V-shape as shown in FIGS. 4B and 4C or may be otherwise configured. As in FIGS. 1 and 2 , the sealable endograft system 1000 of FIG. 3 also has an accommodating proximal cuff 1010 which encloses the terminal lattice structure 1040 as shown and also encloses an adjustment member 1025 that loops through a control device 1030 that is provided to allow increase or decrease in the circumference of the, e.g., circumferential loop of the adjustment member 1025 by the application of rotational torque to the distal aspect of the adjustment tool 1035 emerging from the control device 1030 . The progression of FIGS. 4A to 4C shows the endograft in a relatively contracted form in FIG. 4A , in a partially expanded form in FIG. 4B , and in a fully expanded form in FIG. 4C . As the lattice interruption 1045 is closed in FIGS. 3 and 44 , it can be seen only in FIGS. 4B and 4C . One exemplary configuration for the lattice interruption 1045 can be a woven material that is stretched in the expanded state and attached to the lattice 1041 and, when allowed to reduce, the woven material resist buckling. This configuration allows the diameter to increase beyond the maximum diameter that the graft will allow with the stent alone. [0077] FIG. 5A shows an exemplary embodiment of the control device 1030 in the form of a microcylinder locking mechanism 1050 . This locking mechanism 1050 is changed from a locked state to an unlocked state by an adjustment tool 1060 , which comprises a tool sheath 1062 having a keyed collar portion 1065 . The adjustment tool 1060 is fixed, in both the longitudinal and radial extents, to the remote adjustment tool 1035 . The progression of FIGS. 5A to 5C show how the locking mechanism 1050 is changed from the locked state (in which adjustment of the adjustment member 1025 is prohibited) to the unlocked state (in which adjustment of the adjustment member 1025 is permitted), and, then, back to the locked state. [0078] Before explaining the change between states, the configuration of an exemplary embodiment of the locking mechanism 1050 is described further. The exterior of the locking mechanism 1050 is comprised of a microcylinder 1052 having a set of circumferentially spaced-apart, interior striations 1055 . The locking mechanism 1050 is longitudinally and rotationally fixed to the proximal cuff 1010 . A guide bullet 1070 is received within the hollow, internally striated microcylinder 1052 . The guide bullet 1070 has a longitudinal threaded bore that received therein (in a threaded manner) the adjustment member 1025 . The adjustment member 1025 completely traverses the bore of the guide bullet 1070 and terminates distally of the guide bullet 1070 in a keyed block 1075 that is rotationally fixed to the adjustment member 1025 . The guide bullet 1070 has at least two opposing, flexible tines 1072 that extend radially outward, in a natural state that, together, has a diameter greater than the internal diameter of the locking microcylinder 1052 (the tines can, as well, be spring loaded outwardly). The tines 1072 have a terminal portion that is shaped to fit within a corresponding shaped of each striation 1055 within the microcylinder 1052 . As such, when the tines 1072 are compressed and the guide bullet 1070 is placed within the microcylinder with the adjustment member 1025 threaded therewithin, the tines 1072 press outwardly against the internal surface of the microcylinder 1052 and, when appropriately rotated therein, the tines 1072 each lock within a respective opposing one of the striations 1055 . In such a state, the tines 1072 both form-fittingly and force-fittingly lock within inner striations 1055 when unconstrained. If, for example, there were three tines 1072 separated by 120 degrees each, then the tines 1072 would each lock within a respective one of the striations 1055 that are, also, 120 degrees apart along the interior surface of the microcylinder 1052 . The frictional force of the tines 1072 against the inside surface of the microcylinder 1052 is sufficiently strong to prevent longitudinal movement of the guide bullet 1070 , even if the keyed block 1075 is rotated unless the tines 1072 are removed from their locked position against the interior surface of the microcylinder. In such a configuration, the microcylinder 1052 and the guide bullet 1070 prevent rotation of the adjustment member 1025 without, not only a particular external force applied thereto, but also a removal of the tines 1072 from the interior surface of the microcylinder 1052 . [0079] Rotation of the adjustment member 1025 , therefore, is carried out with the adjustment tool 1060 . The adjustment tool 1060 provides both the ability to rotate the keyed block 1075 but also the ability to separate the tines 1072 from the interior surface of the microcylinder 1052 . To carry out these functions, the tool sheath 1062 has a sufficient cylindrical length to slide between the tines 1072 and the interior surface of the microcylinder 1052 anywhere the tines 1072 are contacting the interior surface. As such, the longitudinal length of the tool sheath 1062 can be, but does not necessarily have to be, as long as the microcylinder 1052 . FIG. 5A shows the microcylinder 1052 with the guide bullet 1070 in a locked position, prior to interface by the remote adjustment tool 1060 . When the adjustment tool 1060 is slid into the microcylinder 1052 , as shown in the progression of FIGS. 5A to 5B , the smooth interior surface of the tool sheath 1062 first slides along the outer surface of the tines and, then, along and past the distal ends of the tines 1072 , at which time the tines 1072 no longer contact the interior surface of the microcylinder 1052 . The orientation of the microcylinder locking mechanism 1050 and the adjustment tool 1060 in FIG. 5B now allows for repositioning of the adjustment member 1025 and relocation of the guide bullet 1070 within the microcylinder 1052 . [0080] The keyed collar portion 1065 has a distal taper 1067 that reduces the outer diameter of the tool sheath 1062 inwards to such an extent that it acts as a funnel to direct the keyed block 1075 directly into the radial center of the keyed collar portion 1065 . At the proximal-most end of the collar portion 1065 is an internal key 1069 having an internal circumferential shape corresponding to an external circumferential shape of the keyed block 1075 . As such, when the adjustment tool 1060 is inserted into the microcylinder 1052 and releases the tines 1072 from the interior surface thereof, the tool sheath 1062 can pass the tines 1072 (wherever they may be inside the microcylinder 1052 ) sufficiently far to permit the keyed block 1075 to slide along the interior distal taper 1067 and press against the internal bore of the key 1069 . With slight rotation either way of the adjustment tool 1060 (by rotation of the adjustment tool 1035 ), the keyed block 1075 will fall into the internal bore of the key 1069 in a form-fit, thereby enabling rotation of the adjustment member 1025 (via keyed block 1075 ) in a corresponding manner to any rotation of the adjustment tool 1035 by a user. [0081] The locking mechanism 1050 is longitudinally and rotationally fixed to the circumferential assembly 1020 such that rotation of the locking mechanism 1050 in a first direction causes a contraction of the circumferential assembly 1020 and rotation of the locking mechanism 1050 in the opposition direction causes an expansion of the circumferential assembly 1020 . As can be seen in FIGS. 5B and 5C , the keyed block 1075 is rotated to cause the guide bullet 1070 to advance towards the keyed block 1075 . FIG. 5C shows the microcylinder locking mechanism 1050 with the adjustment tool 1060 after adjustment and disengagement of the microcylinder locking mechanism 1050 by the adjustment tool 1060 with a fixed repositioning of the guide bullet 1070 and a distal lengthening of the adjustment member 1025 with respect to the microcylinder 1052 . As the final position of the keyed block 1075 is further away from the microcylinder 1052 , and because the microcylinder 1052 is fixed to the control device 1030 of the circumferential assembly 1020 , this exemplary movement of the adjustment member 1025 indicates that the circumferential assembly 1020 has reduced in diameter. [0082] Various alternative embodiments of this locking mechanism are envisioned where a number of the individual parts are fixed or moving with respect to other ones of the parts of the circumferential assembly 1020 , the control device 1030 , the locking mechanism 1050 , and/or the adjustment tool 1060 . In one alternative embodiment of the microcylinder locking mechanism 1050 , the collar portion 1065 of the remote adjustment tool 1060 can contains inner striations (similar to or different from the striations 1055 of the microcylinder 1052 ) that allow it to capture and turn the guide bullet 1070 through removable fixation of the tines 1072 therein (see FIG. 6E ). In such a configuration, the guide bullet 1070 can be fixed rotationally to the adjustment member 1025 . [0083] The inner striations 1055 of the microcylinder 1052 may be grooves, threads, detents, slots, or other surface features sufficient to allow capture of the tines 1072 upon their release as shown in further detail, for example, in the cross-sections of FIGS. 6A to 6G . FIG. 6A is a cross-section along section line A-A of the microcylinder 1052 and guide bullet 1070 of FIG. 5A , in which the tines 1072 having an exemplary triagonal cross-sectional shape are caught within two striations 1055 having an exemplary rectangular cross-sectional shape. FIG. 6B is a cross-section along section line B-B of the tool sheath 1062 of FIG. 5A and illustrates the relatively smooth outer surface of the tool sheath 1062 . FIG. 6C is a cross-section along section line C-C of the microcylinder 1052 of FIG. 5B without the adjustment member 1025 depicted. FIG. 6D is a cross-section along section line D-D of the microcylinder 1052 , the guide bullet 1070 , and the tool sheath 1062 of FIG. 5B , in which the tool sheath 1062 captures the guide bullet 1070 and collapses the tines 1072 , thereby removing the tines 1072 from the striations 1055 of the microcylinder 1052 . [0084] FIG. 6E shows a cross-sectional view of a variation of another exemplary embodiment of the locking mechanism 1050 ′ with the adjustment tool sheath 1062 ′ also having striations 1055 ′ with an exemplary rectangular cross-sectional shape. The tines 1072 are illustrated as expanded within two opposing striations 1055 ′ of the tool sheath 1062 ′. As the tool sheath 1062 ′ has a smooth exterior, the tool sheath 1062 ′ can rotate without friction within the microcylinder 1052 ′. [0085] FIGS. 6F and 6G show cross-sectional views of yet another variation of an exemplary embodiment of the microcylinder locking mechanism 1050 ″ and adjustment tool 1060 ″. The locking mechanism 1050 ″ has a microcylinder 1052 ″ with striations 1055 ″ having an exemplary triangular cross-sectional shape. The adjustment tool sheath 1062 ″ has a smooth exterior and interior to slide within the microcylinder 1052 ″ and to slidably capture the tines 1072 ′″, respectively. The tines 1072 ″ are illustrated as expanded within two opposing triangular striations 1055 ″ of the microcylinder 1052 ″ in FIG. 6F and are captured within the tool sheath 1062 ″ in FIG. 6G . [0086] FIGS. 7A and 7B show longitudinal cross-sectional details of one exemplary embodiment of a locking mechanism 1110 for the adjustment tool 1035 according to the present invention. FIG. 7A shows a locking mechanism 1110 comprising a controllable catch 1115 in a disengaged stated. FIG. 6B shows the locking mechanism 1110 with the controllable catch mechanism 1115 engaged. Once the adjustment member catch 1120 is within the target range 1117 of the locking mechanism, the user can engage a non-illustrated catch deployment device to capture the adjustment member catch 1120 . [0087] FIGS. 8A to 8C show details of still another embodiment of a microcylinder locking mechanism 1150 according to the present invention, in which internal locking tines 1152 , 1154 of unequal length are employed to prevent back rotation from torque buildup upon detachment of the remote adjustment tool 1060 . FIG. 8A shows the locking mechanism 1150 comprised of a microcylinder 1151 and a guide bullet 1153 with internal locking tines 1152 , 1154 of unequal length and an associated adjustment tool 1160 having a tool sheath 1164 prior to engagement of the microcylinder locking mechanism 1150 by the tool sheath 1164 . FIG. 8B shows the tool sheath 1164 of FIG. 8A engaged with the microcylinder locking mechanism 1150 to deflect the tines 1152 , 1154 away from the interior surface of the microcylinder 1151 . FIG. 8C shows the microcylinder locking mechanism 1150 in a locking position different from FIG. 8A after adjustment has occurred and the tool sheath 1164 has been disengaged from the microcylinder 1151 . [0088] FIGS. 9A and 9B show two aspects of details of sheathable retention tines 1130 and a compressible foam sealing gasket 1140 for the proximal terminal aspect of some exemplary embodiments of endografts according to the present invention. FIG. 9A is an axial cross section showing sheathable retention tines 1130 sheathed by an expanded compressible foam gasket 1040 in an exemplary proximal aspect of a sealable endograft system 1000 according to the present invention. FIG. 9B is a perspective view showing sheathable retention tines 1130 exposed and deployed through the compressible foam sealing gasket 1140 disposed at an expanded proximal cuff 1010 in an exemplary endograft according to the present invention. In some exemplary embodiments of the present invention, the direct pressure of the adjustment member 1025 on the footplate 1145 of the tines may be used to extend the sheathable tines 1130 through the compressible foam gasket 1040 and into the wall of a recipient blood vessel. In yet other exemplary embodiments of the present invention, direct pressure of the adjustment member 1025 may exert force on non-illustrated footplate bands that may be attached to or adjacent the footplates 1145 of the tines 1130 and may be used to extend the sheathable tines 1130 through the compressible foam gasket 1040 and into the wall of a recipient blood vessel. Such footplate bands may, themselves, be the base of the sheathable tines 1130 in certain exemplary embodiments of the present invention. Not shown in FIGS. 9A and 9B , the adjustment member 1025 may course though eyelets, other brackets or may otherwise be moveably connected to the footplates 1145 to maintain equal pressure and desired orientation upon expansion of the adjustment member loop. [0089] In the various embodiments of sealable endograft systems according to the present invention, the distal attachment of the endograft to the aortic wall distal to the aneurysm sac may be accomplished in a conventional manner using an expandable lattice component at the distal cuffs, or variations on the adjustable, sealable mechanism disclosed herein may be employed to secure distal seals. The distal seals are subject to lower pressure demands, and the anatomic constraints of sufficient aortic neck distally are generally less problematic than for the proximal seal. [0090] FIGS. 10 to 13 provide anatomic views of another exemplary embodiment of an endograft implant according to the present invention in which the implant is a universal proximal cuff endovascular implant for treatment of an abdominal aortic aneurysm. Endografts with the features shown in the various embodiments of the present invention have unique abilities to accommodate to anatomic variations that would preclude or compromise use of conventional endograft systems. The universal proximal cuff implants of the present invention allow an operator to make use of their ability to securely seal and attach in anatomic sites where conventional endografts cannot be securely placed, and then allow a conventional endograft to securely dock with the universal proximal cuff endovascular implants distally. [0091] Universal proximal cuff endovascular implants of the present invention may be provided with any of the elements disclosed in the present and the incorporated co-pending applications referenced herein. Such elements include, but are not limited to, attachment of radio-opaque monitoring clip assemblies on the outer surfaces of endografts to allow post-implantation monitoring of slippage or endoleak formation by plain radiographs, steerable delivery systems to permit delivery and seal of an endograft in an anatomically angulated or irregular site, and/or auto-accommodation for post-implantation aortic remodeling, [0092] FIG. 10A is an axial cross-sectional view of an exemplary endovascular universal interface cuff 1155 of the present invention to be implanted into an aorta having an aneurysm sac 1170 and an aortic wall 1175 . The universal endovascular interface cuff 1155 has been positioned over an endovascular guidewire 1160 to a desired recipient site A-A′ proximal to the aortic aneurysm sac 1170 . The endovascular universal interface cuff 1155 further comprises an accommodating proximal cuff 1010 and a rigid distal cuff 1200 . FIG. 10B provides a transverse cross-sectional view of the exemplary endovascular interface cuff 1155 of FIG. 10A at the level of A-A′ in FIG. 10A . In FIGS. 10A and 10B , the compressible foam gasket 1140 is uncompressed and, therefore, covers the retention tines 1165 . [0093] In the exemplary embodiment shown in FIG. 10B , the adjustment member 1025 courses in a circumferential loop through eyelets 1180 attached to a series of compression footplates 1185 . The compression footplates 1185 , among other functions, serve to maintain an orientation of the expanding circumferential loop 1035 in a plane transverse to the aortic lumen 1190 , and present a broader pressure contact with the underlying aortic wall 1175 when the circumferential assembly is expanded. The compression footplates 1185 may abut, be attached to, or be contiguous with the retention tines 1165 , which are displaced through the compressed compressible foam gasket 1140 and allowed to enter the aortic wall 1175 for overall device stabilization and retention. While four retention tines 1165 and footplates 1185 are shown, this embodiment is merely exemplary and can be any number. [0094] FIG. 11A shows the same axial cross-sectional view of the endovascular universal interface cuff 1155 of FIG. 10A but after the universal endovascular interface cuff 1155 has expanded to achieve a seal in the aortic wall 1175 . Due to the expansion of the cuff, the foam gasket 1140 becomes compressed, allowing the retention tines 1165 to protrude radially outward to engage the aortic wall 1175 in the desired recipient site A-A′ proximal to the aortic aneurysm sac 1170 . In the exemplary embodiment shown in FIG. 11B , the adjustment member 1025 has expanded to move the eyelets 1180 attached to the footplates 1185 outwards. As is evident, the interior lumen of the circumferential assembly 1020 shown in FIG. 11B has increased substantially as compared to the state shown in FIG. 10B . In FIG. 11B , the compression of the foam gasket 1140 and the engagement of the aortic wall 1175 by the retention tines 1165 creates a firm seal between the universal endovascular interface cuff 1155 and the aortic wall 1175 . [0095] FIG. 12 shows the same axial cross-sectional axial of the universal endovascular interface cuff 1155 of the present invention as in FIGS. 10A and 11A but with delivery of a conventional endograft 1300 into the aortic wall 1175 , which endograft 1300 has been secured within the rigid distal cuff 1200 of the universal endovascular interface cuff 1155 . The endograft 1300 can include an expandable lattice 1310 . FIG. 13 shows the same cross-sectional axial view of an exemplary universal endovascular interface cuff 1155 of the present invention as FIG. 12 but after removal of the endovascular guidewire 1160 and detachment and removal of the adjustment member 1025 . Such removal and detachment can be carried out by a release mechanism 1037 . The distal attachment of the conventional endograft is not shown in FIGS. 12 and 13 , but can be accomplished in the usual manner for conventional endograft implantation sufficient to prevent backfill of the aneurysm sac 1170 from the distal aorta or the iliac vessels. [0096] As shown in FIGS. 10A, 11A, 12, and 13 , the rigid distal cuff 1200 includes, at its exterior, exemplary radio-opaque monitoring clip assemblies 1225 to allow post-implantation monitoring of slippage or endoleak formation and/or auto-accommodation for post-implantation aortic remodeling. Likewise, the rigid distal cuff 1200 can be provided with interior graft retention tines 1227 that add to securing, without leaks, the endograft 1300 to the interior of the rigid distal cuff 1200 . [0097] The tubular endograft body 1005 , the proximal cuff 1010 , the rigid distal cuffs 1200 , and the endograft body 1300 as described herein may be constructed of solid, woven, non-woven, or mesh materials such as, but not limited to, natural or synthetic rubbers, nylon, GORE-TEX®, elastomers, polyisoprenes, polyphosphazenes, polyurethanes, vinyl plastisols, acrylic polyesters, polyvinylpyrrolidone-polyurethane interpolymers, butadiene rubbers, styrene-butadiene rubbers, rubber lattices, DACRON®, PTFE, malleable metals, other biologically compatible materials or a combination of such biologically compatible materials in a molded, woven, or non-woven configuration, coated, non-coated, and other polymers or materials with suitable resilience and pliability qualities. In certain exemplary embodiments according to the present invention, it is desirable for the non-elastic tubular member 1015 and corresponding structures to be pliable to allow for folding or compressibility without allowing elasticity. In certain exemplary embodiments according to the present invention, it is desirable for the accommodating proximal cuff 1010 and corresponding structures to have plasticity and be compressible or foldable. In any given exemplary embodiment, the non-elastic tubular implant body 1015 , the endograft body 1300 , the accommodating proximal cuff 1010 , and corresponding structures may be constructed of the same material of varying elasticity, or these structures may be constructed of different, but compatible materials. [0098] The adjustment members 1025 , the retention tines 1130 , 1165 , and the microcylinders 1030 and other mechanical components as disclosed herein and in all other embodiments of the present invention may be fabricated of any suitably strong biocompatible material, including, but not limited to titanium, stainless steel, cobalt chromium alloys, other metals, other metal alloys, nitinol, plastics, or ceramics. Similarly, the adjustment members 1025 , the retention tines 1130 , 1165 , and the microcylinders 1030 and other mechanical components may be milled, laser cut, lathed, molded, or extruded. [0099] The compressible foam gaskets 1140 as disclosed herein may be any biocompatible foam material of either an open or closed cell structure with sufficient compressibility and resilience to allow rapid recovery in a non-compressed state. In various exemplary embodiments according to the present invention, such foam materials may be viscoelastic foam with a compressible cellular material that has both elastic (spring-like) and viscous (time-dependent) properties. Viscoelastic foam differs from regular foam by having time-dependent behaviors such as creep, stress relaxation, and hysteresis. [0100] FIGS. 14A and 14B show an alternate exemplary embodiment of a sealable endograft system 2000 according to the present invention in two different states. In the view of FIG. 14A , a hinged lattice structure 2100 is attached to an internal or external surface of at least the proximal portion 2210 of an endograft body 2200 (the “lattice” in these figures is only diagrammatic and is not intended to imply that the only possible number of rings of lattice is greater than one). Either the lattice structure 2100 or the endograft body 2200 can be provided with radially displaced retention tines 2105 that, in a non-distended state of the proximal portion 2210 , can be covered within a compressible foam gasket 2300 . In the embodiment shown in FIG. 14A , the distal portion 2220 of the endograft body 2200 comprises a non-distensible material and the proximal portion 2210 of the endograft body 2200 is an accommodating cuff comprising a distensible material forming the proximally terminal aspect of the sealable endograft system 2000 and enclosing the terminal hinged lattice structure 2100 therewithin. [0101] A control system 2400 or jack screw shown in FIGS. 14A and 14B is provided to expand and contract the lattice structure 2100 . In particular, a torque wire 2410 can be fixed at two points 2420 , 2430 longitudinally separate from one another on the lattice structure 2100 . This torque wire 2410 has exterior threads that correspond to threaded bores of one of the two points 2420 , 2430 . Accordingly, when the torque wire 2410 is rotated, the two points 2420 , 2430 of the lattice either approach one another (to expand the proximal portion 2210 ) or retreat from one another (to contract the proximal portion 2210 ) this imparts motion to all contiguously interconnected lattice elements. It is preferred to have the proximal end point 2430 be bored for rotation but fixed longitudinally. In this case, a smooth-bored collar 2440 is fixed to the wall of the graft 2200 , for example, on an interior surface distal of the lattice structure 2100 . When the adjustment tool 1035 is rotated, the torque wire 2410 correspondingly rotates to expand or contract the proximal portion 2210 of the endograft 2200 . In this manner, in comparison to self-expanding prior art stent structures (e.g., made of nitinol) passively open to their greatest extent when relieved from radially inward compression, the lattice structure of the present invention is able to actively open according to the desire of the user surgeon implanting the prosthesis. As such, the opening performed by prior art self-expanding stent structures in endograft prosthesis are referred to herein as “passive opening” or “passive expansion”. In contrast thereto, the expansion performed by the inventive controllable, hinged, lattice structure of the present invention for the disclosed endograft prostheses is referred to herein as “active control” or “active expansion” because it can be actively controlled in both the expansion and contraction directions according to the desire of the user. This is further in contrast to expansion of stent structures using balloon, which case is referred to as “balloon opening” or “balloon expansion” because it occurs only in one direction (expansion) without any ability to contract actively. The single embodiment of the jack screw shown in FIGS. 14A and 14B can be replicated any number of times about the circumference of the lattice structure 2100 [0102] In a non-illustrated alternative to the configuration of the system shown in FIG. 14B , the configuration shown in FIGS. 10A to 11B can be incorporated into the system of FIGS. 14A and 14B to create a hybrid system. The circumferential assembly 1020 can be positioned at the proximal end of the endograft and action of the circumferential loop 1035 within the proximal cuff 1010 , can be used to expand and contract the latticework 2100 . [0103] FIG. 15A is a lateral view of an exemplary embodiment of an adjustable vascular cannula 1230 according to the present invention. As shown in FIG. 15A , such an adjustable vascular cannula 1230 is a generally tubular structure with external cannula walls 1235 defining a cannula lumen 1240 , and comprises a port end 1245 , a cannula body 1250 , and a cannula tip 1255 . As further shown in FIG. 15A , the cannula body 1250 is further provided with a delivery recess 1260 in its external wall structure at or near the junction of the cannula tip 1255 . Further still, the adjustable vascular cannula 1230 of FIG. 15A comprises an adjustable seal device 1265 attached to an adjustment member 1025 such as a torque wire that extends beyond the port end 1245 of the adjustable vascular cannula 1230 as shown in FIG. 15B . The adjustment member 1025 may course through the cannula lumen 1240 , or it may course through an accessory lumen (not shown in FIGS. 15A or 15B ) within the cannula wall 1235 substantially parallel to the cannula lumen 1240 , or it may course externally to the adjustable vascular cannula 1230 as shown partially within and partially outside the lumen 1240 in FIG. 15B . When in a non-deployed state, as shown in FIG. 15B , the adjustable seal device 1265 is substantially flush with the outer diameter of the cannula walls 1235 within the delivery recess 1260 of the cannula body 1250 . [0104] FIG. 15C shows the adjustable seal device 1265 in a deployed state, which is the result of torque applied externally to the adjustment member 1025 by a user. As shown in FIG. 15C , the adjustable seal device 1265 further comprises a hinged adjustable latticework 1270 covered by a sealing cuff 1275 which is constructed of a distensible material. The adjustment member 1025 terminates, for example, in a circumferential loop 1035 within the sealing cuff 1275 , where it may be further covered by a compressible foam gasket 1140 . The adjustment member 1025 may further pass through a locking mechanism 1050 as disclosed elsewhere herein which serves to regulate the torque applied to the circumferential loop 1035 . The hinged adjustable latticework 1270 may further be provided with one or more retention tines 1130 , 1165 , which are radially displaced from the terminal aspect of the hinged adjustable latticework 1270 , and which are enclosed within and covered by the compressible foam gasket 1140 when the adjustable seal device 1265 is not distended. When torque is applied to the adjustment member 1025 by a user, the diameter of the circumferential loop 1035 is increased, displacing the hinged adjustable latticework 1270 as shown in FIG. 15C until the compressible foam gasket 1140 and the sealing cuff 1275 is able to firmly engage the inner wall 1190 of a recipient blood vessel 1175 . A slight additional amount of torque applied to the adjustment member 1025 is, then, sufficient to compress the compressible foam gasket 1140 and allow the retention tines 1130 , 1165 to engage the wall 1190 of the recipient blood vessel 1175 , thus preventing slippage of the cannula during use. In various exemplary embodiments of the present invention, the retention tines 1130 , 1165 may be provided to engage the vessel wall 1190 in a substantially straight manner or at angles varying from about 1 degree to about 179 degrees. The retention tines 1130 , 1165 may be angled axially or longitudinally in various embodiments according to the present invention. After the use of the cannula is completed, the torque of the adjustment member 1025 may be reversed, collapsing the adjustable seal device 1165 , and allowing the compressible foam gasket 1140 to re-expand, thus withdrawing the retention tines 1165 from the vessel wall 1175 and covering the retention tines 1165 to allow atraumatic cannula withdrawal. [0105] Although the foregoing embodiments of the present invention have been described in some detail by way of illustration and example for purposes of clarity and understanding, it will be apparent to those skilled in the art that certain changes and modifications may be practiced within the spirit and scope of the present invention. Therefore, the description and examples presented herein should not be construed to limit the scope of the present invention, the features of which are set forth in the appended claims. [0106] The foregoing description and accompanying drawings illustrate the principles, exemplary embodiments, and modes of operation of the invention. However, the invention should not be construed as being limited to the particular embodiments discussed above. Additional variations of the embodiments discussed above will be appreciated by those skilled in the art and the above-described embodiments should be regarded as illustrative rather than restrictive. Accordingly, it should be appreciated that variations to those embodiments can be made by those skilled in the art without departing from the scope of the invention as defined by the following claims.

Sealable and repositionable implant devices are provided with one or more improvements that increase the ability of implants such as endovascular grafts to be precisely deployed or re-deployed, with better in situ accommodation to the local anatomy of the targeted recipient anatomic site, and/or with the ability for post-deployment adjustment to accommodate anatomic changes that might compromise the efficacy of the implant. A surgical implant includes an implant body and a selectively adjustable assembly attached to the implant body, having adjustable elements, and operable to cause a configuration change in a portion of the implant body and, thereby, permit implantation of the implant body within an anatomic orifice to effect a seal therein under normal physiological conditions.

CROSS-REFERENCE TO RELATED APPLICATIONS [0001] This application is a continuation-in-part of U.S. patent application Ser. No. 14/857,061 filed Sep. 17, 2015, which is hereby incorporated by reference herein. FIELD OF THE DISCLOSURE [0002] The present disclosure relates to neutralization of explosive materials contained in explosives and pyrotechnics. In particular, the disclosure relates to devices and methods for rendering pyrotechnics and ammunition inert or less effective. BACKGROUND [0003] The current worldwide political climate has produced many terrorist and anti-establishment factions that are motivated to create explosive devices from commonly available consumer products. For example, roadside or improvised explosive devices known as IEDs have been encountered in Afghanistan and in Iraq by the U.S. military and in Boston by local police. [0004] A common practice used in constructing an IED involves the acquisition and disassembly of easily acquired consumer grade explosive products such as fireworks or small arms ammunition. The products are disassembled, yielding explosive material, e.g., black powder or other incendiary material. The explosive material is then combined with projectiles such as nails or broken glass and encased in a rigid container such as an aluminum cooking pot. The results are easily concealed and a deadly combination that is both inexpensive and effective. [0005] Consumer grade explosive products contain various explosive materials. For example, gunpowder is a very common chemical explosive and comes in two basic forms, modern, smokeless gunpowder and traditional, black powder gunpowder. Black powder is a mixture of sulfur, charcoal, and potassium nitrate (saltpeter). The sulfur and charcoal act as fuels, and the saltpeter is an oxidizer. The standard composition for gunpowder is about 75% potassium nitrate, about 15% charcoal, and about 10% sulfur (proportions by weight). The ratios can be altered somewhat depending on the purpose of the powder. For instance, power grades of gunpowder, unsuitable for use in firearms but adequate for blasting rock in quarrying operations, have proportions of about 70% nitrate, about 14% charcoal, and about 16% sulfur. Some blasting powder may be made with cheaper sodium nitrate substituted for potassium nitrate and proportions may be as low as about 40% nitrate, about 30% charcoal, and about 30% sulfur. [0006] Most pyrotechnic compositions and explosive materials can be neutralized when mixed with an appropriate combination of inert materials, slowing the burn rate of the explosive material to a non-explosive level that effectively neutralizes the explosive material and renders the explosive material useless for an improvised explosive device. [0007] The prior art addresses the neutralization of explosive devices. However, none of the prior art devices or methods is completely satisfactory in neutralizing explosive materials in consumer products. [0008] For example, U.S. Pat. No. 7,690,287 to Maegerlein, et al. provides a neutralizing assembly for inhibiting operation of an explosive device. The neutralizing assembly will interrupt the function of the explosive device only when the explosive device is misused. The neutralizing assembly includes an interior chamber with a rupturable barrier containing disabling material. The rupturable barrier separates the disabling material from the explosive material. Upon misuse of the device, the rupturable barrier breaks and the disabling material is released from the interior chamber to disable the explosive material. [0009] U.S. Pat. No. 3,738,276 to Picard, et al. discloses a halocarbon gel for stabilizing an explosive material during transport. In use, flexible bags are prepared which contain the explosive material mixed with a desensitizing substance. The bags are placed in a protective gel. The gel prevents the desensitizing substance from evaporating through the flexible bags. When the transport is complete, the bags are removed from the gel. Once the bags are removed from the gel, the desensitizing substance evaporates, thus “arming” the explosive material. [0010] U.S. Patent Publication No. 2011/0124945 to Smylie, et al. discloses a cartridge that is adapted to achieve deactivation of an explosive composition. In Smylie, the explosive composition and the chemical deactivating agent are held in separate chambers of the cartridge separated by a wall. Upon activation, the wall is breached and the deactivating agent and the explosive composition are allowed to mix, thereby rendering the explosive composition inert. [0011] It is, therefore, an object of this disclosure to provide a design for and method of manufacture of products which include an undetectable neutralizing agent that automatically and effectively neutralizes an explosive material upon disassembly. SUMMARY OF THE DISCLOSURE [0012] A concealed amalgamated neutralizer (CAN) is disclosed for the prevention of malicious conversion of consumer fireworks, ammunition, and other pyrotechnic products into dangerous explosive devices, such as an IED. [0013] In a preferred embodiment, a method of manufacture is provided whereby neutralizer material is undetectably situated adjacent to explosive material. The neutralizer material is chosen from various combinations of inert materials such as calcium carbonate, silica, or other inert materials combined with complimentary inert bonding and pigmentation chemicals. The neutralizer material is chosen and modified to mimic the physical characteristics (grain size, density, color) of the explosive material so that when placed side by side with the explosive material, the two components are practically indistinguishable and inseparable. [0014] In one embodiment, the neutralizer material may be a combination of pigmented inert granular constituents. In another embodiment, the neutralizer material may be a liquid or viscous slurry in combination with a source binder and capable of drying into a compact solid. [0015] In another embodiment, a cylindrical design is provided, which positions the explosive material adjacent the neutralizer material along a common central axis. The physical position and/or ratio of the neutralizer material relative to the explosive material can vary to change the extent of the neutralization. [0016] In one embodiment, a temporary build container is provided in the form of a “tube within a tube.” A dry granular explosive material is introduced into the interstitial space between the tubes but excluded from the inner tube. A dry granular neutralizer material of similar color, density, size and texture as the explosive material is then introduced in the inner tube. The inner tube is then removed, allowing the explosive material to contact, but not mix with, the neutralizer material at a boundary interface. The resulting solid cylindrical shape is then packed and sealed, preserving the respective positions of the two components and the boundary interface. [0017] In another embodiment, a spherically shaped device is provided. The neutralizer materials and explosive materials may each be hemispherical and placed “side-by-side.” Temporary physical barriers may be used to separate the components, which are removed during manufacture to create a final product. [0018] In another embodiment of the invention using wet materials, a “layered” product is provided fixed to a substrate. [0019] In each case, the neutralizer material is placed in direct physical contact with the explosive material. The neutralizer material is physically indiscernible from the explosive material, and so the boundary interface between the two is very difficult or impossible to distinguish. Upon disassembly of the product, the neutralizer material is physically mixed with the explosive material, resulting in a combined material that is inert and useless as an explosive. BRIEF DESCRIPTION OF THE DRAWINGS [0020] The disclosed embodiments will be described with reference to the accompanying drawings. [0021] FIG. 1A is a schematic diagram of a portion of a pyrotechnic device in accordance with a preferred embodiment of this disclosure. [0022] FIG. 1B is a schematic diagram of a portion of a pyrotechnic device in accordance with a preferred embodiment of this disclosure. [0023] FIG. 2A is an isometric view of a tube within a tube build container. [0024] FIG. 2B is an isometric view of a preferred embodiment in cylindrical form. [0025] FIG. 3A is an isometric view of a cylindrical layered build container. [0026] FIG. 3B is an isometric view of a preferred embodiment in layered form. [0027] FIG. 4A is a section plan view of spherical side by side build container. [0028] FIG. 4B is a section plan view of a preferred embodiment in spherical form. [0029] FIG. 4C is a section plan view of a spherical build container with a preferred embodiment in spherical form. [0030] FIG. 5 is a flow chart of steps required with assembly of a preferred embodiment of this disclosure. [0031] FIG. 6 is a flow chart of steps to build a spherical pyrotechnic device in accordance with a preferred embodiment of this disclosure. [0032] FIG. 7 is a flow chart of steps to build a spherical pyrotechnic device in accordance with a preferred embodiment of this disclosure. [0033] FIG. 8A is a section plan view of an alternate embodiment resulting from liquid materials. [0034] FIG. 8B is a section plan view of an alternate embodiment resulting from liquid materials as it is being made. [0035] FIG. 9 is a flow chart of steps required with assembly of a preferred embodiment of this disclosure. [0036] FIG. 10 is a section plan view of an article of manufacture including a preferred embodiment of this disclosure. [0037] FIG. 11 is a flow chart of steps for assembly of an article of manufacture including a preferred embodiment of this disclosure. [0038] FIG. 12 is a section plan view of a Roman candle in accordance with a preferred embodiment of this disclosure. [0039] FIG. 13 is a flow chart of steps to build a Roman candle in accordance with a preferred embodiment of this disclosure. [0040] FIG. 14 is an isometric view of a pyrotechnic assembly in accordance with a preferred embodiment of this disclosure. [0041] FIG. 15 is a flow chart of steps to build a pyrotechnic assembly in accordance with a preferred embodiment of this disclosure. [0042] FIG. 16 is an isometric view of a pyrotechnic assembly in accordance with a preferred embodiment of this disclosure. [0043] FIG. 17 is a flow chart of steps to roll a pyrotechnic device in accordance with a preferred embodiment of this disclosure. [0044] FIG. 18 is a detail view of a pyrotechnic device in accordance with a preferred embodiment of this disclosure. [0045] FIG. 19 is a flow chart of steps to build a device using a shell case in accordance with a preferred embodiment of this disclosure. [0046] FIG. 20 is a cross section view of a pyrotechnic pigeon in accordance with a preferred embodiment of this disclosure. [0047] FIG. 21A is a flow chart of steps to build a pyrotechnic pigeon in accordance with a preferred embodiment of this disclosure. [0048] FIGS. 21B to 211 are cross section views of a pyrotechnic pigeon as it is being built in accordance with a preferred embodiment of this disclosure. DETAILED DESCRIPTION [0049] Referring to FIG. 1A , portion 100 of a pyrotechnic or explosive device is shown that includes concealed amalgamated neutralizer 104 to prevent the use of explosive composition 114 in other devices. Portion 100 comprises housing 102 , which acts to enclose and/or support concealed amalgamated neutralizer 104 and explosive composition 114 . Concealed amalgamated neutralizer 104 and explosive composition 114 are positioned with or adjacent to each other. Interface 132 is an indiscernible boundary interface between concealed amalgamated neutralizer 104 and explosive composition 114 and is where concealed amalgamated neutralizer 104 touches explosive composition 114 . Example pyrotechnic devices that comprise portion 100 include ammunition (such as shotgun shell 1000 of FIG. 10 ), fireworks (such as Roman candle 1200 of FIG. 12 ), and other explosive devices (such as a training target comprising the devices of FIGS. 8A, 8B and 18 and percussion caps). [0050] Concealed amalgamated neutralizer 104 is a composition having a color and grain size that is indiscernible from the color and grain size of explosive composition 114 . When mixed sufficiently with explosive composition 114 , explosive power of the resulting mixture is reduced as compared to the explosive power of explosive composition 114 so as to prevent the use of explosive composition 114 outside of housing 102 . Concealed amalgamated neutralizer 104 comprises non-inert material 106 , inert material 108 , and binding agent 112 . Concealed amalgamated neutralizer 104 may be formed from a slurry, such as neutralizer slurry 124 of FIG. 1B . [0051] In alternative embodiments, concealed amalgamated neutralizer 104 is formed without being processed from a neutralizer slurry. As an example, concealed amalgamated neutralizer 104 may be formed from a dry powder. [0052] Materials used as non-inert material 106 include aluminum and may optionally comprise or form a pigment. Non-inert material 106 may include materials similar to fuel 116 of explosive composition 114 . Non-inert material 106 alters the fuel to oxidizer ratio of explosive composition 114 and/or provides different burn characteristics so as to reduce the explosiveness of explosive composition 114 when explosive composition 114 is combined with concealed amalgamated neutralizer 104 outside of housing 102 . [0053] Materials used in inert material 108 include magnesium silicate and chalk and may optionally comprise or form a pigment. Inert material 108 does not burn or explode and acts to reduce the explosiveness of explosive composition 114 when explosive composition 114 is combined with concealed amalgamated neutralizer 104 outside of housing 102 . [0054] Materials used as binding agent 112 of concealed amalgamated neutralizer 104 include cellulose and shellac and also include materials similar to materials used as binding agent 122 of explosive composition 114 . Binding agent 112 acts to bind the components of concealed amalgamated neutralizer 104 together and prevent the components of concealed amalgamated neutralizer 104 from mixing with explosive composition 114 while concealed amalgamated neutralizer 104 and explosive composition 114 are contained within the pyrotechnic device comprising portion 100 . [0055] Referring to FIG. 1B , a substrate 103 may also be used to support various embodiments where a liquid binder is necessary. Neutralizer slurry 124 and explosive slurry 128 are formed on top of substrate 103 . Interface 133 is an indiscernible boundary interface between neutralizer slurry 124 and explosive slurry 128 . Neutralizer slurry 124 and explosive slurry 128 are positioned with or adjacent to each other and touch each other at interface 133 . [0056] Neutralizer slurry 124 is used to form concealed amalgamated neutralizer 104 . Neutralizer slurry 124 includes non-inert material 106 , inert material 108 , and binding agent 112 . Neutralizer slurry 124 also includes solvent 126 . Once positioned with respect to substrate 103 , neutralizer slurry 124 is allowed to solidify by withdrawal of solvent 126 , e.g., via vaporization, to form concealed amalgamated neutralizer 104 as a solid or to give concealed amalgamated neutralizer 104 a more solid-like form. [0057] Materials used as solvent 126 include methyl ethyl ketone (MEK), cellulose thinners, isopropanol, alcohol, water, hydrogen peroxide, liquefied petroleum gas (LPG), and liquid nitrogen. Solvent 126 dissolves the other components of neutralizer slurry 124 and allows neutralizer slurry 124 to be processed in a more liquid-like fashion as compared to concealed amalgamated neutralizer 104 . [0058] Explosive composition 114 is an explosive material, also known as a pyrotechnic composition, comprising one or more fuels 116 , oxidizers 118 , and additives 120 , and binding agents 122 . Fuels 116 and oxidizers 118 interact chemically to release energy, additives 120 add additional properties, and binding agents 122 bind explosive composition 114 together. Explosive composition 114 is formed from explosive slurry 128 . [0059] In alternative embodiments, explosive composition 114 is formed without being processed from explosive slurry 128 . As an example, explosive composition 114 may be formed from a dry powder. [0060] Materials used as fuel 116 include: metals, metal hydrides, metal carbides, metalloids, non-metallic inorganics, carbon based compounds, organic chemicals, and organic polymers and resins. Metal fuels include: aluminum, magnesium, magnalium, iron, steel, zirconium, titanium, ferrotitanium, ferrosilicon, manganese, zinc, copper, brass, tungsten, zirconium-nickel alloy. Metal hydride fuels include: titanium(II) hydride, zirconium(II) hydride, aluminum hydride, and decaborane. Metal carbides used as fuels include zirconium carbide. Metalloids used as fuels include: silicon, boron, and antimony. Non-metallic inorganic fuels include: sulfur, red phosphorus, white phosphorus, calcium silicide, antimony trisulfide, arsenic sulfide (realgar), phosphorus trisulfide, calcium phosphide, and potassium thiocyanate. Carbon based fuels include: carbon, charcoal, graphite, carbon black, asphaltum, and wood flour. Organic chemical fuels include: sodium benzoate, sodium salicylate, gallic acid, potassium picrate, terephthalic acid, hexamine, anthracene, naphthalene, lactose, dextrose, sucrose, sorbitol, dextrin, stearin, stearic acid, and hexachloroethane. Organic polymer and resin fuels include: fluoropolymers (such as Teflon and Viton), hydroxyl-terminated polybutadiene (HTPB), carboxyl-terminated polybutadiene (CTPB), polybutadiene acrylonitrile (PBAN), polysulfide, polyurethane, polyisobutylene, nitrocellulose, polyethylene, polyvinyl chloride, polyvinylidene chloride, shellac, and accroides resin (red gum). [0061] Materials used as oxidizers 118 include: perchlorates, chlorates, nitrates, permanganates, chromates, oxides and peroxides, sulfates, organic chemicals, and others. Perchlorate oxidizers include: potassium perchlorate, ammonium perchlorate, and nitronium perchlorate. Chlorate oxidizers include: potassium chlorate, barium chlorate, and sodium chlorate. Nitrates include: potassium nitrate, sodium nitrate, calcium nitrate, ammonium nitrate, barium nitrate, strontium nitrate, and cesium nitrate. Permanganate oxidizers include: potassium permanganate and ammonium permanganate. Chromate oxidizers include: barium chromate, lead chromate, and potassium dichromate. Oxide and peroxide oxidizers include: barium peroxide, strontium peroxide, lead tetroxide, lead dioxide, bismuth trioxide, iron(II) oxide, iron(III) oxide, manganese(IV) oxide, chromium(III) oxide, and tin(IV) oxide. Sulfate oxidizers include: barium sulfate, calcium sulfate, potassium sulfate, sodium sulfate, and strontium sulfate. Organic oxidizers include: guanidine nitrate, hexanitroethane, cyclotrimethylene trinitramine, and cyclotetramethylene tetranitramine. Other oxidizers include: sulfur, Teflon, and boron. [0062] Materials used as additives 120 include materials that act as: coolants, flame suppressants, opacifiers, colorants, chlorine donors, catalysts, stabilizers, anticaking agents, plasticizers, curing and crosslinking agents, and bonding agents. Coolants include: diatomaceous earth, alumina, silica, magnesium oxide, carbonates including strontium carbonate, and oximide. Flame suppressants include: potassium nitrate and potassium sulfate. Opacifiers include carbon black and graphite. Colorants include: salts of metals (including barium, strontium, calcium, sodium, and copper), copper metal, and copper acetoarsenite with potassium perchlorate. Chlorine donors include: polyvinyl chloride, polyvinylidene chloride, vinylidene chloride, chlorinated paraffins, chlorinated rubber, hexachloroethane, hexachlorobenzene, and other organochlorides and inorganic chlorides (e.g., ammonium chloride, mercurous chloride), as well as perchlorates and chlorates. Catalysts include: ammonium dichromate, iron(III) oxide, hydrated ferric oxide, manganese dioxide, potassium dichromate, copper chromite, lead salicylate, lead stearate, lead 2-ethylhexoate, copper salicylate, copper stearate, lithium fluoride, n-butyl ferrocene, di-n-butyl ferrocene. Stabilizers include: carbonates (e.g., sodium, calcium, or barium carbonate), alkaline materials, boric acid, organic nitrated amines (such as 2-nitrodiphenylamine), petroleum jelly, castor oil, linseed oil, ethyl centralite, and 2-nitrodiphenylamine. Anticaking agents include: fumed silica, graphite, and magnesium carbonate. Plasticizers: include dioctyl adipate, isodecyl pelargonate, and dioctyl phthalate as well as other energetic materials such as: nitroglycerine, butanetriol trinitrate, dinitrotoluene, trimethylolethane trinitrate, diethylene glycol dinitrate, triethylene glycol dinitrate, bis(2,2-dinitropropyl)formal, bis(2,2-dinitropropyl)acetal, 2,2,2-trinitroethyl 2-nitroxyethyl ether, and others. Curing and crosslinking agents include: paraquinone dioxime, toluene-2, 4-diisocyanate, tris(1-(2-methyl) aziridinyl) phosphine oxide, N,N,O-tri(1,2-epoxy propyl)-4-aminophenol, and isophorone diisocyanate. Bonding agents include tris(1-(2-methyl) azirinidyl) phosphine oxide and triethanolamine. [0063] Materials used as binding agents 122 include: gums, resins and polymers, such as: acacia gum, red gum, guar gum, copal, cellulose, carboxymethyl cellulose, nitrocellulose, rice starch, cornstarch, shellac, dextrin, hydroxyl-terminated polybutadiene (HTPB), polybutadiene acrylonitrile (PBAN), polyethylene, and polyvinyl chloride (PVC). [0064] Explosive slurry 128 is used to form explosive composition 114 . Explosive slurry 128 includes fuel 116 , oxidizer 118 , additives 120 , and binding agent 122 . Explosive slurry 128 also includes solvent 130 . Once positioned with respect to housing 102 , explosive slurry 128 is allowed to solidify by withdrawal of solvent 130 , e.g., via vaporization, to form explosive slurry 128 as a solid or to give explosive slurry 128 more solid-like form. [0065] Materials used as solvent 130 include methyl ethyl ketone (MEK), cellulose thinners, isopropanol, alcohol, water, and hydrogen peroxide. Solvent 130 dissolves the other components of explosive slurry 128 and allows explosive slurry 128 to be processed in a more liquid-like fashion as compared to explosive composition 114 . [0066] Table 1 below shows typical components of dry granular explosive materials, dry neutralizer materials, coloring agents, and ratios required to neutralize the explosive materials in several preferred embodiments. The ratios indicated are by weight, but similar ratios may also be made by volume. The percentage composition of the explosive materials can vary by as much as plus or minus 15%. The percentage composition of the neutralizer materials can vary by as much as plus or minus 15%. The composition ratios can vary by as much as plus or minus 25%. [0000] TABLE 1 Dry Explosive Dry Neutralizer Coloring DEM:DIM Materials Materials Agents (by weight) 70% potassium 65% magnesium Aluminum 3:2 chlorate silicate 30% aluminum 30% aluminum 5% ackroyd resin 75% potassium nitrate Silica Carbon slurry 3:1 15% charcoal 10% sulfur 70% potassium nitrate Silica Carbon slurry 3:1 14% charcoal 16% sulfur 40% sodium nitrate Chalk Carbon black 3:2 30% charcoal 30% sulfur 75% potassium nitrate Barium Lamp black 6:5 19% carbon 6% sulfur [0067] Table 2 below shows typical components of explosive materials, neutralizer materials, pigmentation, solvents, and ratios. The percentage composition of the explosive materials can vary by as much as plus or minus 15%. The percentage composition of the neutralizer materials can vary by as much as plus or minus 15%. The composition ratios can vary by as much as plus or minus 25%. [0000] TABLE 2 Explosive Neutralizer EM:IM:Sol Materials Materials Pigmentation Solvents (by weight) 75% potassium Silica Carbon black Alcohol 3:1:1 nitrate 15% charcoal 10% sulfur 70% potassium Chalk Lamp black Water 3:2:2 nitrate 14% charcoal 16% sulfur 40% sodium Barium Aluminum Isopropanol 6:5:4 nitrate pigment 30% charcoal (ultramarine) 30% sulfur 75% potassium Saw dust Vine black Liquid 11:9:9 nitrate nitrogen 19% carbon 6% sulfur [0068] Tables 3-5 below show typical components of neutralizers, solvents, pigments, and explosive compounds, any of which may be used in pyrotechnic devices in accordance with this disclosure. Table 3 below includes a list of neutralizers and solvents, any of which may be used in pyrotechnic devices. [0000] TABLE 3 Neutralizers Solvents Talcum Methyl ethyl ketone (MEK) Chaulk Cellulose thinners Barrium Isopropanol Manganese Water Aluminum Alcohol Silica Hydrogen peroxide Saw dust Liquefied petroleum gas Calcium carbonate Liquid nitrogen Barite Potters clay [0069] Table 4 below shows a list of pigments, any of which may be used in pyrotechnic devices. A pigment that is used in portion 100 of pyrotechnic device may form part of non-inert material 106 or part of inert material 108 , depending on the chemical composition of the pigment. When a pigment is used to tint concealed amalgamated neutralizer 104 , a sufficient amount is used to coat and color the granules formed from non-inert material 106 and inert material 108 within concealed amalgamated neutralizer 104 . The amount or proportion of pigment may vary depending on the grain size of the granules formed from non-inert material 106 and inert material 108 within concealed amalgamated neutralizer 104 . The pigment may be introduced to concealed amalgamated neutralizer 104 in the form of a dye. Similarly, the granules of the inert materials may be washed with a pigment or dye for a time sufficient to change their color to approximate the color of the granules of the non-inert material. The grainsize of the pigmented inert material can be controlled by sifting with an appropriate wire mesh or other method as known in the art. The mesh size is chosen to approximate the size of the non-inert material. [0000] TABLE 4 Pigments Aluminum pigments: ultramarine violet, ultramarine Antimony pigments: antimony white Arsenic pigments: orpiment natural monoclinic arsenic sulfide (As 2 S 3 ) Barium pigments: barium sulfate Biological pigments: alizarin, alizarin crimson, gamboge, cochineal red, rose madder, indigo, Indian yellow, Tyrian purple Cadmium pigments: cadmium yellow, cadmium red, cadmium green, cadmium orange, cadmium sulfoselenide (CdSe) Carbon pigments: carbon black, ivory black (bone char), vine black, lamp black, India ink Chromium pigments: chrome green, viridian, chrome yellow, chrome orange Clay earth pigments (iron oxides): yellow ochre, raw sienna, burnt sienna, raw umber, burnt umber Cobalt pigments: cobalt violet, cobalt blue, cerulean blue, aureolin (cobalt yellow) Copper pigments: Azurite, Han purple, Han blue, Egyptian blue, Malachite, Paris green, Scheele's Green, Phthalocyanine Blue BN, Phthalocyanine Green G, verdigris, viridian Iron pigments: Prussian blue, yellow ochre, iron black Iron oxide pigments: sanguine, caput mortuum, oxide red, red ochre, Venetian red, burnt sienna Lead pigments: lead white, cremnitz white, Naples yellow, red lead Manganese pigments: manganese violet Mercury pigments: vermilion Organic pigments: quinacridone, magenta, phthalo green, phthalo blue, pigment red 170, diarylide yellow Tin pigments: mosaic gold Titanium pigments: titanium yellow, titanium beige, titanium white, titanium black Ultramarine pigments: ultramarine, ultramarine green shade Zinc pigments: zinc white, zinc ferrite India ink [0070] Table 5 below shows typical explosive compounds, any of which may be used in pyrotechnic devices in accordance with this disclosure. Table 5 includes the following acronyms (among others): trinitrotoluene (TNT), ammonium nitrate (AN), ammonium nitrate fuel oil (ANFO), triethylenetetramine (TETA), nitromethane (NM), penthrite (PETN), research department explosive (RDX), erythritol tetranitrate (ETN), high-velocity military explosive (HMX), polyurethane (PU), polycaprolactone (PCP), trimethylolethane trinitrate (TMETN), hydroxyl-terminated polybutadiene (HTPB), alkyl acrylate copolymer (ACM), dioctyl adipate (DOA), ammonium perchlorate (AP), nitrocellulose (NC), and isopropyl nitrate (IPN). [0000] TABLE 5 Explosive compounds Aluminum powder (30%) + Potassium chlorate (70%) Amatol (50% TNT + 50% AN) Amatol (80% TNT + 20% AN) Ammonium nitrate (AN + <0.5% H 2 O) ANFO (94% AN + 6% fuel oil) ANNMAL (66% AN + 25% NM + 5% Al + 3% C + 1% TETA) Black powder (75% KNO 3 + 19% C + 6% S) Blasting powder Chopin's Composition (10% PETN + 15% RDX + 72% ETN) Composition A-5 (98% RDX + 2% stearic acid) Composition B (63% RDX + 36% TNT + 1% wax) Composition C-3 (78% RDX) Composition C-4 (91% RDX) DADNE (1,1-diamino-2,2-dinitroethene, FOX-7) DDF (4,4′-Dinitro-3,3′-diazenofuroxan) Diethylene glycol dinitrate (DEGDN) Dinitrobenzene (DNB) Erythritol tetranitrate (ETN) Ethylene glycol dinitrate (EGDN) Flash powder Gelatine (92% NG + 7% nitrocellulose) Heptanitrocubane (HNC) Hexamine dinitrate (HDN) Hexanitrobenzene (HNB) Hexanitrostilbene (HNS) Hexogen (RDX) HMTD (hexamine peroxide) HNIW (CL-20) Hydrazine mononitrate Hydromite ® 600 (AN water emulsion) MEDINA (Methylene dinitroamine) Mixture: 24% nitrobenzene + 76% TNM Mixture: 30% nitrobenzene + 70% nitrogen tetroxide Nitrocellulose (13.5% N, NC) Nitroglycerin (NG) Nitroguanidine Nitromethane (NM) Nitrourea Nobel's Dynamite (75% NG + 23% diatomite) Nitrotriazolon (NTO) Octanitrocubane (ONC) Octogen (HMX grade B) Octol (80% HMX + 19% TNT + 1% DNT) PBXIH-135 EB (42% HMX, 33% Al, 25% PCP-TMETN's system) PBXN-109 (64% RDX, 20% Al, 16% HTPB's system) PBW-11 (96% HMX, 1% ACM, 3% DOA) PBW-126 (22% NTO, 20% RDX, 20% AP, 26% Al, 12% PU's system) Penthrite (PETN) Pentolite (56% PETN + 44% TNT) Picric acid (TNP) Plastics Gel ® (45% PETN + 45% NG + 5% DEGDN + 4% NC) RISAL P (50% IPN + 28% RDX + 15% Al + 4% Mg + 1% Zr + 2% NC) Semtex 1A (76% PETN + 6% RDX) Tanerit Simply ® (93% granulated AN + 6% red P + 1% C) acetone peroxide (TATP) Tetryl Tetrytol (70% tetryl + 30% TNT) trinitroazetidine (TNAZ) Torpex (aka HBX, 41% RDX + 40% TNT + 18% Al + 1% wax) Triaminotrinitrobenzene (TATB) Trinitrobenzene (TNB) Trinitrotoluene (TNT) Tritonal (80% TNT + 20% aluminium) [0071] Referring to FIG. 2A , build container 202 is shown. Build container 202 is a generally hollow cylinder having sidewall 204 , open end 206 , and closed end 208 defining interior space 205 . In one embodiment, number 20 cardboard is used to form the ends and walls. Other structural materials such as mylar or vinyl will suffice. Build container 202 is used in a preferred method of assembling generally cylindrical shaped devices containing various combinations of dry compositions of explosive and neutralizer materials, as will be further described. Inner tube 210 is removably affixed within the interior of build container 202 by means common in the art, such as a suitably releasable adhesive. In the preferred embodiment, inner tube 210 is located co-axially with build container 202 , however inner tube 210 may be positioned anywhere within interior 205 . Although a single inner tube is depicted within build container 202 , it will be understood that a plurality of inner tubes may be installed inside build container 202 . Inner tube 210 has an exterior cylindrical shaped surface 212 and an open end 214 defining interior space 215 . Neutralizer material is loaded into interior space 215 , which is inside of interior space 205 , and the explosive material is loaded into interior space 205 outside of interior space 215 . Those skilled in the art will understand that shapes other than cylindrical may be used for inner tube 210 and/or build container 202 such as elliptical, rectangular, and triangular. It is further understood that the size of inner tube 210 relative to build container 202 can be changed depending on the ratio of neutralizer material to explosive material required to properly render the explosive material useless. Additionally, the overall volume of the assembled device may vary depending on intended use of the device. [0072] It should be understood that the positions of the explosive and neutralizer materials could be reversed so that explosive material is loaded into interior space 215 , which is inside of interior space 205 , and the neutralizer material is loaded into interior space 205 outside of interior space 215 . Furthermore, the relative dimensions of the build container and the inner tube organize functions of the ratio of explosive and neutralizer materials. [0073] FIG. 2B shows an assembled device 222 containing neutralizer material 220 and explosive material 230 separated by a boundary interface 225 . Neutralizer material 220 is comprised of components that match explosive material 230 such that neutralizer material 220 is indiscernible from explosive material 230 . Neutralizer material 220 is chosen to approximate the grain size and color of explosive material 230 . Boundary interface 225 is where explosive material 230 contacts neutralizer material 220 within assembled device 222 . Since neutralizer material 220 is indiscernible from explosive material 230 , boundary interface 225 is not visible. [0074] Referring to FIG. 3A , alternate build container 302 is shown. Build container 302 is a generally hollow cylinder having sidewall 304 , open end 306 , and closed end 308 defining interior space 305 . Build container 302 is used for assembling generally disc shaped, layered devices. [0075] FIG. 3B shows an assembled device 322 made from build container 302 in which dry manufacture neutralizer material 320 is layered on top of explosive material 330 . In an alternate embodiment, explosive material 330 is layered on top of neutralizer material 320 . Explosive material 330 is separated from neutralizer material 320 by boundary interface 325 . [0076] FIG. 4A shows an alternate build container 402 . Build container 402 is comprised of two hollow, semi-spherical halves 404 and 406 . Half 404 defines interior space 408 and half 406 defines interior space 410 . A disk shaped separation barrier 409 may be affixed to either half 404 or 406 to contain the explosive material and neutralizer material during assembly. [0077] FIG. 4B shows an assembled device 422 made from build container 402 . Explosive material 430 is separated from neutralizer material 420 by boundary interface 425 . Boundary interface 425 is imperceptible upon visual inspection. [0078] In an alternate spherical arrangement shown in FIG. 4C , build container 402 is used to create a spherical shaped device comprised of a spherical core surrounded by a larger sphere. Explosive material 430 is a hollow sphere shape including a spherical shaped core of neutralizer material 420 . It should be understood by those skilled in the art that an arrangement of neutralizer material surrounding explosive material would be equally effective. Imperceptible boundary interface 426 is provided between explosive material 430 and neutralizer material 420 . [0079] For simplicity in FIGS. 1-4 , detonators, primers, fuses, igniters, casings, plugs, etc. are not shown as each device may require different combinations of these elements typically found in various consumer fireworks, ammunition, and other pyrotechnic products. Some devices use other sources of ignition such as heat or impact. [0080] Referring to FIG. 5 , the steps involved with constructing a device using generally dry materials are shown. At step 502 , an explosive material is chosen. The proper explosive material will be chosen based on its intended use. At step 504 the grain size of the explosive material is identified. If the explosive material contains multiple components each having different grains sizes, each grain size will be identified. At step 506 , the color of the explosive material is identified. At step 508 , a matching neutralizer material with the identified grain size is chosen. The neutralizer material and the level of neutralization desired are chosen according to Table 1 for dry materials or Table 2 for slurries. At step 510 , if the color of the neutralizer material does not match the explosive material, then the neutralizer material is colored using a pigment or dye to match the explosive material. In a different embodiment, a charcoal dye is employed to tint the neutralizer material. At step 512 , the explosive material is introduced into a build container. At step 514 , the neutralizer material is introduced into the build container, and if necessary, the build container is assembled. If necessary, at step 516 , the materials introduced in the build container are compacted. At step 518 , the separation barrier is removed from the build container. At step 520 , any ancillary components required for the device, such as plugs, primers, fuses, detonators, etc., are installed and the assembled device is wrapped in appropriate casing. [0081] Referring to FIG. 6 , one or more steps involved with constructing a spherical pyrotechnic device using generally inert materials are shown. At step 602 , an explosive material is chosen. The proper explosive material will be chosen based on its intended use. At step 604 , the dry density of the explosive material is identified. At step 606 , the color of the dried explosive material is identified. At step 608 , a slurry is prepared from the explosive material and the appropriate solvent or liquid. At step 610 , the neutralizer material with the identified dry density is chosen. At step 612 , a neutralizer slurry is prepared using the neutralizer material and proper pigmentation and solvent. [0082] At step 614 , the neutralizer slurry is rolled into a sphere. In a preferred embodiment, the neutralizer slurry is rolled into a sphere through the use of a scoop. In one preferred embodiment, a scoop is used which is part number ZEROLL 1020 available from Centinal Restaurant Products of Indianapolis, Ind. [0083] At step 616 , the neutralizer slurry is optionally allowed to at least partially solidify so that the sphere of the neutralizer slurry will maintain its geometry during subsequent processing. At step 618 , the explosive slurry is rolled into a sphere such that the volume of the sphere of the neutralizer slurry and the volume of the sphere of the explosive slurry forms a selected ratio, e.g., 2:3 or about 40% to about 60%. [0084] At step 620 , the sphere of neutralizer slurry is implanted into the sphere of the explosive slurry. The sphere of neutralizer slurry is implanted into substantially the center of the sphere of the explosive slurry to create a substantially uniform spherical explosive profile. In other embodiments, the shape and position of the neutralizer slurry within the sphere of explosive slurry is selected to create a non-uniform explosive profile that is not spherical. [0085] At step 622 , the volume of explosive slurry into which the sphere of neutralizer slurry was implanted is rolled again to reform a spherical shape. At step 624 , the explosive slurry is allowed to solidify and, if it is not already solidified, the neutralizer slurry within the sphere of explosive slurry is also optionally allowed to solidify and dry. The sphere comprising the solidified explosive slurry and the neutralizer slurry may then be used to form a pyrotechnic device. [0086] Referring to FIG. 7 , one or more steps involved with constructing a preferred device is shown. At step 702 , an explosive material is chosen. The proper explosive material will be chosen based on its intended use. At step 704 , the dry density of the explosive material is identified. At step 706 , the color of the dried explosive material is identified. At step 708 , a slurry is prepared from the explosive material and the appropriate solvent or liquid. At step 710 , the neutralizer material with the identified dry density is chosen. At step 712 , a neutralizer slurry is prepared using the neutralizer material and proper pigmentation and solvent. At step 714 , the neutralizer slurry is rolled into a sphere. At step 716 , the neutralizer slurry is optionally allowed to at least partially solidify so that the sphere of the neutralizer slurry will maintain its geometry during subsequent processing. At step 718 , explosive slurry is applied and rolled onto the sphere of partially solidified neutralizer slurry. At step 720 , the explosive slurry is allowed to solidify and, if it is not already solidified, the neutralizer slurry within the sphere of explosive slurry is also optionally allowed to solidify and dry. The sphere comprising the solidified explosive slurry and the neutralizer slurry may then be used to form a pyrotechnical device. [0087] FIG. 8A shows an alternate embodiment of device 824 constructed on substrate 840 . Substrate 840 is preferably paper, but may also take the form of other planar surfaces or objects. Explosive material 830 is adhered to substrate 840 . Neutralizer material 820 is adhered to both explosive material 830 and substrate 840 thereby encapsulating the explosive material and forming boundary interface 826 . Device 824 is manufactured from slurry compositions of explosive materials and neutralizer materials as will be further described. [0088] The thickness of explosive material 830 on substrate 840 is substantially uniform along the surface of substrate 840 , except at the outer edges. The thickness of neutralizer material 820 on explosive material 830 and on substrate 840 is also substantially uniform, except at the outer edges. In alternative embodiments, the thicknesses may vary. For example, when device 824 embodies a target training dummy, a thickness of explosive material 830 at substantially the center of the target training dummy may be increased and a thickness of neutralizer material 820 may be reduced to retain a similar overall thickness. In this manner, a different pyrotechnic and visual effect is achieved so that a hit substantially in the center of the target training dummy is distinguishable from a hit that is not substantially in the center of the target training dummy. [0089] FIG. 8B shows an alternate embodiment of device 824 as a layer of neutralizer material 820 is being applied to explosive material 830 . Neutralizer material 820 is prepared in tank or hopper 852 and then applied to explosive material 830 on substrate 840 . Tank or hopper 852 includes an outlet 854 and a valve 856 at the underside of tank or hopper 852 , and outlet 854 is controlled by a valve 856 . The valve 856 can be adjusted to control the volume of the neutralizer slurry dispensed. One of the tank or hopper 852 or the substrate 840 is moved in a direction so that a controlled amount of neutralizer material 820 is applied to explosive material 830 . In a preferred embodiment, the thickness of neutralizer material 820 is substantially the same as the thickness of explosive material 830 . In alternative embodiments, the thicknesses of neutralizer material 820 and explosive material 830 may vary. [0090] Referring to FIG. 9 , the steps involved with constructing a preferred device is shown. At step 932 , an explosive material is chosen. The proper explosive material will be chosen based on its intended use. At step 934 , the dry density of the explosive material is identified. At step 936 , the color of the dried explosive material is identified. At step 937 , a slurry is prepared from the explosive material and the appropriate solvent or liquid. At step 938 , the neutralizer material with the identified dry density is chosen. The neutralizer material is selected from Table 3. [0091] At step 940 , a neutralizer slurry is prepared using the neutralizer material, proper pigmentation and solvent. In a preferred embodiment, the neutralizer slurry is an embodiment of neutralizer slurry 124 of FIG. 1B and is prepared by placing all of the ingredients or components of neutralizer slurry into a tank or hopper in which the ingredients or components are mixed. [0092] At step 942 , the explosive slurry is applied to the substrate. At step 944 , the explosive slurry is allowed to solidify and dry. [0093] At step 946 , the neutralizer slurry is applied to the dried explosive slurry and the substrate. In a preferred embodiment, the underside of a tank or hopper, such as tank or hopper 852 of FIG. 8B , in which the neutralizer slurry was prepared includes an outlet, such as outlet 854 , controlled by a valve, such as valve 856 . The valve can be adjusted to control the volume of the neutralizer slurry dispensed. The valve is placed over the article on which neutralizer slurry 820 is to be applied. For example, the article may comprise substrate 840 and explosive material 830 of FIGS. 8A and 8B . After placement of the valve, the valve is actuated to dispense a selected amount of the neutralizer slurry onto the article to achieve a desired ratio between the amount of neutralizer slurry and the amount of explosive slurry on the article. [0094] At step 948 , the neutralizer slurry is allowed to solidify and dry. [0095] In one preferred embodiment, an article of manufacture, in this case a shotgun shell, is produced according to this disclosure. Referring to FIG. 10 , an article of manufacture, shotgun shell 1000 , is shown. Shotgun shell 1000 includes casing 1002 enclosed on one end by base 1004 . Primer 1006 extends through base 1004 and is positioned adjacent generally cylindrically shaped concealed amalgamated device 1008 . Concealed amalgamated device 1008 is comprised of neutralizer material 1010 separated from explosive material 1012 by boundary interface 1014 . Adjacent the explosive material and neutralizer material is wad 1016 . Shot 1018 is shown adjacent wad 1016 . Crimped closure 1017 is shown opposite base 1004 . [0096] Referring to FIG. 11 , a flowchart showing the steps involved in loading a shotgun shell casing incorporating a preferred embodiment of the device. At step 1104 , the primer is pressed into the base. A separation barrier in the form of a cylindrical Mylar tube is placed in the casing adjacent the base at step 1106 . In a preferred embodiment, the tube is located coaxially with the primer. At step 1108 , gunpowder is loaded into the casing within the interior of the separation barrier. At step 1109 , the neutralizer material is chosen to match the color and grain size of the gunpowder. Choice of the neutralizer material includes the optional selection of a pigment or dye used to match the color of the neutralizer material to the color of the gunpowder. At step 1110 , the neutralizer material is loaded into the casing surrounding the separation barrier. At step 1112 , the separation barrier is removed. At step 1114 , a wad is loaded and pressed within the casing. At step 1116 , shot is loaded and pressed into the casing. At step 1118 , the casing is crimped closed. [0097] In use, should the shotgun shell be disassembled, the neutralizer material is automatically and undetectably mixed with the explosive material. Since the neutralizer material cannot be easily separated from the explosive material, the mixture effectively cannot be used to form an improvised explosive device. [0098] In one preferred embodiment, an article of manufacture, in this case a pyrotechnic device commonly referred to as a Roman candle, is produced according to this disclosure. Referring to FIG. 12 , an article of manufacture, Roman candle 1200 , is shown. Roman candle 1200 includes one or more: fuse 1202 , delay charges 1204 and 1212 , stars 1206 and 1214 , lift charges 1208 and 1216 , neutralizer rings 1210 and 1218 , clay plug 1220 , and paper wrapping 1222 . [0099] Fuse 1202 is connected to a first delay charge 1204 . Fuse 1202 is a burning fuse that, when lit, burns for a selected amount of time based on the length of fuse 1202 and where fuse 1202 is lit along the length of fuse 1202 . Fuse 1202 passes fire to and ignites delay charge 1204 . [0100] Delay charge 1204 is connected to fuse 1202 and packed on top of a first star 1206 , lifting charge 1208 , and shaped neutralizer ring 1210 . Delay charge 1204 comprises a pyrotechnic composition that burns at a slow constant rate that is not significantly affected by temperature or pressure and is used to control timing of the pyrotechnic device, i.e., Roman candle 1200 . After being ignited by fuse 1202 , first delay charge 1204 burns for a selected amount of time based on the composition, height, volume, and density of delay charge 1204 , and then ignites one or more of star 1206 and lift charge 1208 . Delay charge 1204 delays the time between the burning of fuse 1202 and ignition of star 1206 and lift charge 1208 . [0101] Star 1206 is positioned between delay charge 1204 and lift charge 1208 . Star 1206 comprises a pyrotechnic composition selected to provide a visual effect, including burning a certain color or creating a spark effect once first star 1206 is ignited. Star 1206 is coated with black powder to aid the ignition of star 1206 and aid the ignition of lift charge 1208 . [0102] First lift charge 1208 is positioned between first delay charge 1204 and second delay charge 1212 and is in contact with first star 1206 and first shaped neutralizer ring 1210 . First lift charge 1208 comprises an explosive material, such as granulated black powder or any compound selected from Table 5, and is used to shoot first star 1206 out of Roman candle 1200 and to ignite second delay charge 1212 . Ignition of first lift charge 1208 causes first star 1206 to shoot out of Roman candle 1200 with a velocity based on one or more of the composition, size, shape, and position of first lift charge 1208 within Roman candle 1200 . As depicted in FIG. 12 , first lift charge 1208 is shaped substantially as an inverted frustum of a right angle cone with a diameter of the base contacting first delay charge 1204 being larger than a diameter of the base contacting second delay charge 1212 . The shape of lift charge 1208 in conjunction with the shape of neutralizer ring 1210 operate to control the blast profile of the explosion created when lift charge 1208 is ignited. The shape of an inverted frustum provides for the explosion created by the ignition of first lift charge 1208 to be directed out through the top of Roman candle 1200 while still allowing for sufficient contact area with second delay charge 1212 to pass fire onto and ignite second delay charge 1212 after first lift charge 1208 is ignited. [0103] Neutralizer ring 1210 surrounds the conically slanted side of lift charge 1208 and is positioned between delay charge 1204 and delay charge 1212 . Neutralizer ring 1210 is a ring of material comprising an inert material that, as described above, is indiscernible from the explosive material of lift charge 1208 and that, if mixed with the explosive material of lift charge 1208 , results in a composition having a substantially reduced explosiveness. Material of shaped neutralizer ring 1210 has a grain size and color matching that of the grain size and color of material of lift charge 1208 so that the interface between shaped neutralizer ring 1210 and lift charge 1208 is indiscernible. [0104] Delay charge 1212 , star 1214 , lift charge 1216 , and neutralizer ring 1218 operate in a similar fashion as delay charge 1204 , star 1206 , lift charge 1208 , and neutralizer ring 1210 , but may have the same or different compositions, sizes, shapes, positions, and geometries and provide for the same or different specific effects. [0105] Clay plug 1220 is a bottom layer of Roman candle 1200 beneath the combination of second lift charge 1216 and neutralizer ring 1218 . Clay plug 1220 prevents fire from second lift charge 1216 from escaping through the bottom of Roman candle 1200 and prevents lift charge 1216 from being ignited from below. [0106] Paper wrapping 1222 surrounds the sides of Roman candle 1200 forming a cylindrical shape. Paper wrapping 1222 protects Roman candle 1200 when not in use and acts as a muzzle to direct stars 1206 and 1214 when they are shot out of the top of Roman candle by lift charges 1208 and 1216 , respectively. [0107] Referring to FIG. 13 , one or more steps involved with constructing a pyrotechnic device commonly referred to as a Roman candle is shown. At step 1302 , an explosive material is chosen. The proper explosive material will be chosen based on its intended use and may be selected from the explosive compounds from Table 5. At step 1304 , the dry density of the explosive material is identified. At step 1306 , the color of the dried explosive material is identified. At step 1308 , the lift charge, star and delay charge are prepared using explosive material. At step 1310 , the neutralizer material with the identified dry density is selected from the neutralizers listed in Table 3. At step 1312 , a neutralizer powder is prepared using the neutralizer material and proper pigmentation and solvent selected from Tables 3-4. [0108] At step 1314 , a paper tube is prepared to receive the clay plug, one or more lift charges, one or more stars, one or more delay charges and neutralizer powder. The paper tube may be placed vertically so that the materials may be introduced from the top of the tube. At step 1316 , a clay plug is inserted into the bottom of tube that directs the explosions from the lift charge out through the top of the tube. At step 1318 , a separation barrier is inserted into the tube. The separation barrier may include a slant to be slightly conical in shape so that the lift charge is formed as a frustum. At step 1320 , the lift charge is inserted into the tube inside the separation barrier, after which one or more stars are placed on top of the lift charge. At step 1322 , neutralizer powder is inserted into the tube outside of the separation barrier. The neutralizer powder has the same grain size and color as the lift charge. At step 1324 , the separation barrier is removed and the interface between the lift charge and the neutralizer is indiscernible due to the selected properties of the neutralizer powder. At step 1326 , a delay charge is inserted into the tube and packed down so that the lift charge, stars, neutralizer powder, and delay charge will not mix during subsequent handling and processing. At step 1328 , steps 1318 - 1326 are repeated for a desired number of stages for the pyrotechnic device. At step 1330 , a fuse is introduced into the tube that contacts the top-most delay charge. [0109] In one preferred embodiment, an article of manufacture, in this case a pyrotechnic assembly, is produced according to this disclosure. Referring then to FIG. 14 , an article of manufacture, pyrotechnic assembly 1400 , is shown. Pyrotechnic assembly 1400 includes: paper 1402 , slurry 1404 , fuse 1406 , and solidified material 1408 . [0110] Paper 1402 forms an outer shell for a pyrotechnic device created from assembling pyrotechnic assembly 1400 . Prior to rolling paper 1402 to form a cylinder, slurry 1404 is placed on paper 1402 , solidified material 1408 is placed onto slurry 1404 , and fuse 1406 is positioned. After positioning slurry 1404 , solidified material 1408 , and fuse 1406 onto paper 1402 , paper 1402 is rolled to form a cylindrical pyrotechnic device. [0111] Slurry 1404 is positioned on paper 1402 between paper 1402 and solidified material 1408 prior to rolling paper 1402 . After rolling, slurry 1404 forms a substantially continuous layer around solidified material 1408 . One of slurry 1404 and solidified material 1408 comprises neutralizer material (e.g., concealed amalgamated neutralizer 104 of FIG. 1A ) and the other of slurry 1404 and solidified material 1408 comprises explosive material (e.g., explosive composition 114 of FIG. 1A ). After solidifying, the boundary between the material of slurry 1404 and the material of solidified material 1408 will be indiscernible upon visual inspection. The volume of slurry 1404 is sufficient so that when the material of slurry 1404 is randomly mixed with the material of solidified material 1408 , the explosiveness of the combined mixed material is substantially reduced. [0112] Fuse 1406 is positioned to pass flame to explosive material comprised by one of slurry 1404 and solidified material 1408 . Fuse 1406 contacts both slurry 1404 and solidified material 1408 so that fuse 1406 contacts both the inert material of one of slurry 1404 and solidified material 1408 and the explosive material of the other of slurry 1404 and solidified material 1408 . By contacting both slurry 1404 and solidified material 1408 , the position of fuse 1406 does not provide an indication of whether solidified material 1408 or slurry 1404 comprises explosive material in the final assembled device. [0113] In an alternative embodiment where solidified material 1408 comprises the explosive material, fuse 1406 may be positioned within and incorporated into solidified material 1408 prior to the solidification of solidified material 1408 . With fuse 1406 incorporated into solidified material 1408 , placement of solidified material 1408 also positions fuse 1406 with respect to paper 1402 of assembly 1400 . [0114] Solidified material 1408 is positioned on slurry 1404 prior to rolling paper 1402 and contacts fuse 1406 . After rolling pyrotechnic assembly 1400 into a pyrotechnic device, solidified material 1408 is located in substantially the center of the pyrotechnic device. In alternative embodiments, solidified material 1408 may be positioned away from the center of the pyrotechnic device and create a different explosion profile as compared to when the solidified material 1408 is placed in the center of the pyrotechnic device. [0115] Referring to FIG. 15 , one or more steps involved with constructing a pyrotechnic device by rolling single portions of explosive material and neutralizer material into a cylinder is shown. At step 1502 , an explosive material is chosen from Table 5. The proper explosive material will be chosen based on its intended use. At step 1504 , the dry density of the explosive material is identified. At step 1506 , the color of the dried explosive material is identified. At step 1508 , an explosive slurry is using the explosive material and the appropriate solvent or liquid. At step 1510 , the neutralizer material with the identified dry density is chosen. At step 1512 , a neutralizer slurry is prepared using the neutralizer material and proper pigmentation and solvent or liquid. [0116] At step 1514 , paper is prepared for creating the pyrotechnic device. The paper is formed as a square or rectangular sheet with appropriate dimensions of thickness, length, and width to form the exterior of the pyrotechnic device. At step 1516 , a first slurry is applied to the paper. The first slurry is one or the other of the explosive slurry and the neutralizer slurry. At step 1518 and prior to introducing the second slurry to the first slurry, the second slurry is allowed to at least partially solidify to form a solidified material or paste that is thicker than the first slurry to aid further processing steps. The second slurry is different from the first slurry and is the other of the explosive slurry or the neutralizer slurry. At step 1520 , the solidified material made from the second slurry is positioned onto the first slurry. [0117] At step 1522 , a fuse is introduced between the solidified material and the first slurry so as to contact the explosive material in one or the other of the first slurry and the second slurry. In alternative embodiments, the fuse is introduced into the second slurry prior to solidification of the second slurry. At step 1524 , the paper is rolled into a cylindrical shape. The process or rolling the paper surrounds the entirety of the solidified material with the first slurry and positions the solidified material substantially in the center of the cylinder created by rolling the paper. Positioning the solidified material in the center of the cylinder gives the pyrotechnic device a substantially uniform blast profile along the circumference of the cylinder. In alternative embodiments, the solidified material is positioned off center so that the pyrotechnic device will not contain a substantially uniform blast profile along the circumference of the cylinder [0118] In one preferred embodiment, an article of manufacture, in this case a pyrotechnic assembly, is produced according to this disclosure. Referring to FIG. 16 , an article of manufacture, assembly 1600 , is shown that forms an embodiment of portion 100 of a pyrotechnic device of FIG. 1A . Assembly 1600 includes: paper 1602 , explosive compound 1604 , and neutralizer compound 1606 . [0119] Paper 1602 is a substrate onto which explosive compound 1604 and neutralizer compound 1606 are applied. After application of explosive compound 1604 and neutralizer compound 1606 onto paper 1602 , paper 1602 is rolled from one end in direction 1608 to form a cylinder. A fuse for igniting explosive compound 1604 may be introduced to assembly 1600 before or after rolling paper 1602 into a cylinder. After assembly into pyrotechnic device, paper 1602 protects the pyrotechnic device from unwanted ignition. [0120] Explosive compound 1604 is any explosive material and is applied to paper 1602 as a paste or slurry to stick between multiple layers of paper 1602 after paper 1602 is rolled. The width of each portion of explosive compound 1604 applied to paper 1602 is substantially uniform. In alternative embodiments, the width of each portion of explosive compound 1604 applied to paper 1602 may vary along the length of paper 1602 . The overall ratio of the volume of explosive compound 1604 to the volume of neutralizer compound 1606 is such that, if explosive compound 1604 and neutralizer compound 1606 are removed from a pyrotechnic device created from assembly 1600 and mixed, then the resulting mixture would have a substantially reduced explosive effectiveness. [0121] Neutralizer compound 1606 is any neutralizer material and is also applied to paper 1602 as a paste or slurry to stick between multiple layers of paper 1602 after paper 1602 is rolled. The width of each portion of neutralizer compound 1606 applied to paper 1602 is substantially uniform and is less than the width of the portions of explosive compound 1604 . When dried, neutralizer compound 1606 has a grain size that substantially matches the grain size of explosive compound 1604 . Neutralizer compound 1606 includes pigmentation so that the color of neutralizer compound 1606 substantially matches the color of explosive compound 1604 . The boundary interface between the portions of explosive compound 1604 and neutralizer compound 1606 are indiscernible upon final assembly due to the matching grain size and color between explosive compound 1604 and neutralizer compound 1606 . [0122] In alternative embodiments, the width of each portion of explosive compound 1604 applied to paper 1602 may vary along the length of paper 1602 . [0123] Referring to FIG. 17 , one or more steps involved with constructing a pyrotechnic device by rolling multiple portions of explosive material and neutralizer material is shown. At step 1702 , an explosive material is chosen from Table 5. The proper explosive material will be chosen based on its intended use. At step 1704 , the dry density of the explosive material is identified. At step 1706 , the color of the dried explosive material is identified. At step 1708 , a slurry is prepared from the explosive material and the appropriate solvent or liquid. At step 1710 , the neutralizer material with the identified dry density is chosen. At step 1712 , a neutralizer slurry is prepared using the neutralizer material and proper pigmentation and solvent. [0124] At step 1714 , paper is prepared as a substrate to receive the explosive slurry and neutralizer slurry. The paper is sliced into a selected length and width suitable for rolling. At step 1716 , explosive slurry and neutralizer slurry are applied to the paper in alternating portions, as shown in FIG. 16 . The width of the portions may be uniform or vary based on the location of the portion with respect to the leading edge of the paper that gets rolled first and the trailing edge of the paper that gets rolled last. For example, portions closer to the trailing edge may have a longer width as compared to portions closer to the leading edge [0125] At step 1718 , the paper with the applied explosive slurry and neutralizer slurry is rolled into a cylindrical shape so that each portion of explosive compound contacts two portions of neutralizer compound and two layers of paper. Similarly, each portion of neutralizer compound contacts two portions of explosive compound and two layers of paper. [0126] At step 1720 , a fuse is inserted into the cylinder created by rolling the paper. The fuse is inserted so as to contact at least one portion of explosive slurry. At step 1722 , at least the explosive slurry is allowed to solidify and optionally the neutralizer is also allowed to solidify. [0127] At step 1720 , the explosive slurry is allowed to solidify as well as the neutralizer slurry. The cylindrically shaped roll comprising the solidified explosive slurry and the neutralizer slurry may then be used to form a pyrotechnical device. With the color, grain size, and dry density being substantially similar, the interfaces between portions of explosive material and neutralizer material in the rolled cylinder are indiscernible upon visual inspection and the explosive material is indistinguishable from the neutralizer material. Removal of the explosive material would also remove the neutralizer material so that attempted use of the explosive material in an improvised explosive device would mix the explosive material with the neutralizer material and reduce the effectiveness of the explosive material in the improvised explosive device. [0128] In one preferred embodiment, an article of manufacture, in this case pyrotechnic device 1800 forms, for example, an instant hit recognition flare or pyrotechnic target, and is produced according to this disclosure. Referring to FIG. 18 , an article of manufacture, pyrotechnic device 1800 , is shown that forms an embodiment of portion 100 of a pyrotechnic device of FIG. 1A . Pyrotechnic device 1800 includes: cardboard lid 1801 , concealed amalgamated neutralizer 1802 , pyrotechnic composition 1803 , imperceptible boundary layer 1804 , and shell case 1805 . [0129] Cardboard lid 1801 and shell case 1805 form an embodiment of housing 102 of FIG. 1A . Cardboard lid 1801 is fitted to the top of shell case 1805 and presses against concealed amalgamated neutralizer 1802 to compact and maintain the shape and position of concealed amalgamated neutralizer 1802 and pyrotechnic composition 1803 within pyrotechnic device 1800 . [0130] Concealed amalgamated neutralizer 1802 is layered on top of pyrotechnic composition 1803 and is held in place by cardboard lid 1801 and shell casing 1805 . Pyrotechnic composition 1803 is an embodiment of explosive composition 114 , is layered on top of shell case floor 1806 , and is held in place by shell casing 1805 . When concealed amalgamated neutralizer 1802 is mixed with pyrotechnic composition 1803 outside of pyrotechnic device 1800 , such as in an improvised explosive device, the explosive power of the resulting mixture is reduced as compared to the explosive power of pyrotechnic composition 1803 . [0131] Imperceptible boundary layer 1804 is present at the interface or junction between concealed amalgamated neutralizer 1802 and pyrotechnic composition 1803 . Concealed amalgamated neutralizer 1802 is selected, processed, and manufactured to comprise a grain shape, grain size, color, and density that substantially matches the grain shape, grain size, color, and density of pyrotechnic composition 1803 so that imperceptible boundary layer 1804 cannot be perceived upon visual inspection. [0132] Shell case 1805 comprises shell case floor 1806 and contains concealed amalgamated neutralizer 1802 and pyrotechnic composition 1803 . Shell case 1805 presses against concealed amalgamated neutralizer 1802 and pyrotechnic composition 1803 to compact and maintain the shape and position of concealed amalgamated neutralizer 1802 and pyrotechnic composition 1803 within pyrotechnic device 1800 . [0133] Referring to FIG. 19 , the steps involved with constructing a pyrotechnic device with concealed amalgamated neutralizer as used in an instant hit recognition flare or pyrotechnic target using a shell case is shown. At step 1902 , an explosive material, also known as a pyrotechnic composition, is chosen. The proper explosive material will be chosen based on its intended use. At step 1904 the grain size of the explosive material is identified. If the explosive material contains multiple components each having different grains sizes, each grain size will be identified. At step 1906 , the color of the explosive material is identified. At step 1908 , a matching neutralizer material, also known as a concealed amalgamated neutralizer or a concealed amalgamated neutralizer component, with the identified grain size is chosen. The neutralizer material and the level of neutralization desired is chosen according to Table 1 for dry materials or Table 2 for slurries. At step 1910 , if the color of the neutralizer material does not match the explosive material, then the neutralizer material is colored to match the explosive material using one or more pigments or dyes. In a different embodiment, a charcoal dye is employed to tint the neutralizer material. At step 1912 , the explosive material is introduced into a shell case. At step 1914 , the neutralizer material is introduced into the shell case, and if necessary, the shell case is assembled. If necessary, at step 1916 , the materials introduced in the build container are compacted. At step 1918 , a cardboard lid is installed onto and fitted to the shell case. In alternative embodiments, the materials are compacted after installation of the cardboard lid instead of or in addition to being compacted prior to installation of the cardboard lid. At step 1920 , any ancillary components required for the device, such as plugs, primers, fuses, detonators, etc., are installed. [0134] In one preferred embodiment, an article of manufacture, in this case a pyrotechnic pigeon, is produced according to this disclosure. Referring to FIG. 20 , an article of manufacture, pyrotechnic pigeon 2000 , is shown that includes an embodiment of portion 100 of a pyrotechnic device of FIG. 1A . Pyrotechnic pigeon 2000 is a target configured for target shooting. Pyrotechnic pigeon 2000 includes substrate layer 2002 , first plastic layer 2004 , first material layer 2006 , second material layer 2008 , and second plastic layer 2010 . The sizes and thicknesses of the layers are not shown to scale. In certain embodiments, pyrotechnic pigeon 2000 comprises a standard clay pigeon to which first plastic layer 2004 , first material layer 2006 , second material layer 2008 , and second plastic layer 2010 are applied. [0135] Substrate layer 2002 includes a step-shaped edge 2012 at the circumference of pyrotechnic pigeon 2000 . Step-shaped edge 2012 allows for pyrotechnic pigeon 2000 to be guided and rotated as it is launched from a clay pigeon launcher. Substrate layer 2002 acts as a substrate upon which is formed first plastic layer 2004 , first material layer 2006 , second material layer 2008 , and second plastic layer 2010 . Substrate layer 2002 contacts one or more layers of plastic material. Substrate layer 2002 comprises any clay, plastic, metal, concrete, limestone, pitch, or other material that is suitable for making a targets for clay pigeon shooting, also known as clay target shooting. [0136] First plastic layer 2004 is positioned between substrate layer 2002 and first material layer 2006 . First plastic layer 2004 protects first material layer 2006 from substrate layer 2002 . First plastic layer 2004 adheres the combination of first plastic layer 2004 , first material layer 2006 , second material layer 2008 , and second plastic layer 2010 to substrate layer 2002 . [0137] First material layer 2006 is positioned between first plastic layer 2004 and second material layer 2008 . Second material layer 2008 is positioned between first material layer 2006 and second plastic layer 2010 . [0138] When first material layer 2006 is the explosive material, second material layer 2008 is the neutralizer material. When first material layer 2006 is the neutralizer material, second material layer 2008 is the explosive material. The neutralizer material is selected and processed to have the same color, density, dry weight, and grain size as the explosive material so that the junction between first material layer 2006 and second material layer 2008 is formed as an indiscernible boundary layer. The ratio of explosive material to neutralizer material is such that, if explosive material and neutralizer material were removed from pyrotechnic pigeon 2000 and mixed, then the resulting mixture would have substantially reduced usefulness as a propellant or explosive, such as in an improvised explosive device. [0139] Second plastic layer 2010 is placed onto second material layer 2008 and substrate layer 2002 . Second plastic layer 2010 surrounds the outer edges of each of first plastic layer 2004 , first material layer 2006 , and second material layer 2008 . Second plastic layer 2010 protects and supports first material layer 2006 and second material layer 2008 . Combined, first plastic layer 2004 and second plastic layer 2010 operate to seal, protect, and encapsulate first material layer 2006 and second material layer 2008 from external moisture and humidity. [0140] First plastic layer 2004 and second plastic layer 2010 may be homogeneous or heterogeneous and comprise any form of plastic, including: acrylic, acrylonitrile butadiene styrene (ABS), diallyl-phthalate (DAP), epoxy resin, high impact polystyrene (HIPS), high-density polyethylene (HDPE), low-density polyethylene (LDPE), medium-density polyethylene (MDPE), melamine resin, phenol formaldehyde resin (PF), polyactic acid (PLA), polyamide (PA) (nylon), polybenzimidazole (PBI), polycarbonate (PC), polycyanurate, polyester (PE), polyether sulfone (PES), polyetherether ketone (PEEK), polyetherimide (PEI), polyethylene (PE), polyethylene terephthalate (PET), polyimide (PI), polymethyl methacrylate (PMMA), polyphenylene oxide (PPO), polyphenylene sulfide (PPS), polypropylene (PP), polystyrene (PS), polytetrafluoroethylene (PTFE), polyurethane (PU), polyvinyl chloride (PVC), polyvinylidene chloride (PVDC), urea-formaldehyde, and vulcanized rubber. In one preferred embodiment, first plastic layer 2004 comprises an acrylic resin and is enhanced for adhesive properties to ensure the combination of first plastic layer 2004 , first material layer 2006 , second material layer 2008 , and second plastic layer 2010 adheres to substrate layer 2002 . Second plastic layer 2010 is enhanced for brittleness to protect the placement and positioning of the combination of first plastic layer 2004 , first material layer 2006 , second material layer 2008 , and second plastic layer 2010 on top of substrate layer 2002 during transport and handling. [0141] Referring to FIGS. 21A to 211 , FIG. 21A is a flow chart depicting steps used to create a pyrotechnic pigeon, such as pyrotechnic pigeon 2000 of FIG. 20 , and FIGS. 21B to 211 are cross section views of a pyrotechnic pigeon as it is being built with the steps of FIG. 21A . [0142] At step 2102 , an explosive material is chosen to be used for the pyrotechnic pigeon. The proper explosive material will be chosen based on its intended use and may be selected from the explosive compounds from Table 5. In one preferred embodiment, explosive material includes black powder and one or more pyrotechnic stars that become visible when the pyrotechnic pigeon is hit. In another preferred embodiment, explosive material includes flash powder to create a visible flash and audible noise when the pyrotechnic pigeon is hit. [0143] At step 2104 , the properties of the explosive material are identified, which include the color, weight, density, and grain size of the explosive material in its final dry form in the pyrotechnic pigeon. [0144] At step 2106 , the explosive material is prepared for processing. In one preferred embodiment, the explosive material is formed as an explosive slurry that can be particlized or sprayed onto a surface. [0145] At step 2108 , a neutralizer material is chosen to be used for the pyrotechnic pigeon. The neutralizer material chosen has similar properties as the explosive material or can be processed to have properties that are substantially similar to the properties of the explosive material. [0146] At step 2110 , the neutralizer material is prepared for processing. If the neutralizer material chosen does not have an appropriate color, then a pigment is added to the neutralizer material that give the neutralizer material a color that is substantially the same as or is indiscernible from the color of the explosive material. In one preferred embodiment, the neutralizer material is formed as a neutralizer slurry that can be particlized or sprayed onto a surface. [0147] At step 2112 , substrate layer 2002 (shown in FIG. 21B ) is formed. In one preferred embodiment, substrate layer 2002 is formed by compacting a mixture of pitch and pulverized limestone in a mold to form the shape of the substrate layer 2002 . In another preferred embodiment, substrate layer 2002 is a pre-manufactured clay pigeon. [0148] At step 2114 , outer guide 2130 (shown in FIG. 21C ) is placed onto substrate layer 2002 . In one preferred embodiment, outer guide 2130 is cylindrically shaped and includes step-shaped edge 2132 that matches a portion of step-shaped edge 2012 of substrate layer 2002 . Matching step-shaped edge 2132 of outer guide 2130 to the portion of step-shaped edge 2012 of substrate layer 2002 centers and seals outer guide 2130 to substrate layer 2002 so that material applied within outer guide 2130 is appropriately placed onto substrate layer 2002 without leaking onto or reaching step-shaped edge 2012 of substrate layer 2002 . In certain embodiments, shapes other than or in addition to a step are used to match or key outer guide 2130 to substrate layer 2002 . [0149] At step 2116 , inner guide 2134 (shown in FIG. 21D ) is placed onto substrate layer 2002 within outer guide 2130 . Inner guide 2134 is cylindrically shaped with an outer circumference that is similar to the inner circumference of outer guide 2130 so that inner guide 2134 fits within outer guide 2130 and is centered with respect to outer guide 2130 and to substrate layer 2002 . A bottom edge of inner guide 2134 contacts a top surface of substrate layer 2002 to prevent material applied within inner guide 2134 from reaching outer guide 2130 on the top surface of substrate layer 2002 . [0150] At step 2118 , first plastic layer 2004 (shown in FIG. 21E ) is formed. In one preferred embodiment, first plastic layer 2004 is sprayed onto substrate layer 2002 within inner guide 2134 . Inner guide 2134 prevents the application of first plastic layer 2004 from reaching the inner edge of outer guide 2130 . [0151] At step 2120 , first material layer 2006 (shown in FIG. 21F ) is formed. In one preferred embodiment, first material layer 2006 is an explosive material that is sprayed onto first plastic layer 2004 within inner guide 2134 . Inner guide 2134 prevents the application of first material layer 2006 from reaching the inner edge of outer guide 2130 . [0152] At step 2122 , second material layer 2008 (shown in FIG. 21G ) is formed. In one preferred embodiment, second material layer 2008 is a neutralizer material that is sprayed onto first material layer 2006 within inner guide 2134 . Inner guide 2134 prevents the application of second material layer 2008 from reaching the inner edge of outer guide 2130 . [0153] At step 2124 , inner guide 2134 is removed (shown in FIG. 21H ). Removing inner guide 2134 exposes outer edges of first plastic layer 2004 , first material layer 2006 , and second material layer 2008 . Removing inner guide 2134 also exposes the portion of the top surface of substrate layer 2002 that was covered by the bottom surface of inner guide 2134 . [0154] At step 2126 , second plastic layer 2010 (shown in FIG. 21I ) is formed. In one preferred embodiment, second plastic layer 2010 is sprayed so that the application of second plastic layer covers second material layer 2008 , reaches the edges of first material layer 2006 and first plastic layer 2004 within outer guide 2130 , and reaches the top surface of substrate layer 2002 that was covered by the bottom surface of inner guide 2134 . Outer guide 2130 prevents the application of second plastic layer 2010 from reaching step-shaped edge 2012 of substrate layer 2002 . [0155] At step 2128 , outer guide 2130 is removed from the fully formed pyrotechnic pigeon, such as pyrotechnic pigeon 2000 (shown in FIG. 20 ). Removing outer guide 2130 exposes the outer edge of second plastic layer 2010 and the portion of the top surface of substrate layer 2002 that was covered by the bottom surface of outer guide 2130 . [0156] It will be appreciated by those skilled in the art that modifications can be made to the embodiments disclosed and remain within the inventive concept. Therefore, this invention is not limited to the specific embodiments disclosed, but is intended to cover changes within the scope and spirit of the claims.

A concealed amalgamated neutralizer covertly combines neutralizer material comprised of various combinations of inert materials such as calcium carbonate or silicates with common explosive material for the prevention of malicious use of the explosive material in improvised explosive devices. The concealed amalgamated neutralizer device may vary in shape, size, and color and is therefore adaptable to varying methods of containment typified by common pyrotechnic products. The neutralizer material mimics the explosive material of the pyrotechnic products without detection. Upon disassembly of a concealed amalgamated neutralizer device, the neutralizer material is mixed with and neutralizes the explosive material rendering the explosive material useless as a component for an improvised explosive device.

CROSS-REFERENCE TO RELATED APPLICATIONS This application is a continuation-in-part of U.S. patent application Ser. No. 13/275,871, which claims priority to Korean Patent Application No. 10-2010-0107014, filed on Oct. 29, 2010, and all the benefits accruing therefrom under 35 U.S.C. §119, the content of which in its entirety is herein incorporated by reference. INCORPORATION-BY-REFERENCE OF MATERIAL SUBMITTED ELECTRONICALLY Incorporated by reference in its entirety herein is a computer-readable nucleotide/amino acid sequence listing submitted concurrently herewith and identified as follows: One 588 Byte ASCII (Text) file named “710372SequenceListing.txt,” created on Jun. 8, 2012. BACKGROUND 1. Field The present disclosure relates to a device for separating and/or disrupting cells, and more particularly, to a micro-device for separating and/or disrupting cells and a method for disrupting cells using the micro-device. 2. Description of the Related Art A nucleic acid-based micro total analysis system (“μTAS”) is a promising platform for analyzing molecules. In such an analytic system, all analysis processes, e.g., a nucleic acid sample preparation process, an amplification process and a detection process, are integrated on a microchip and are automated. Real-time polymerase chain reaction (“PCR”) process has been widely used since additional post-PCR analysis processes, e.g., electrophoresis and fluorescent imaging, may be omitted, thereby saving time and reducing complexity. However, the sample preparation process for obtaining a nucleic acid suitable for PCR from a raw sample has not been effectively integrated into the entire analysis process on a microchip. Impurities in a sample may directly affect the specificity and sensitivity of PCR. Accordingly, elimination of impurities in a sample may substantially improve the amplification process. Moreover, a target analyte in a sample is concentrated from an initial large volume to a small volume, and thus the initial sample having a volume of milliliters (ml) may be processed in the micro-device. Such a property improves PCR sensitivity, and is of great merit in using the micro-device beyond the capacities of laboratory equipment. Accordingly, inclusion of such functions in the device for nucleic acid preparation may substantially improve the application of the microfluidic analysis system. The nucleic acid preparation process typically includes cell disruption for releasing nucleic acid from inside a cell to outside the cell. Cell disruption may include a process for disrupting cell membranes or cell walls. In particular, gram positive bacteria have a very thick peptidoglycan layer. Accordingly, it is more difficult to disrupt the cell membranes of gram positive bacteria than those of gram negative bacteria. Staphylococcus aureus (“ S. aureus ”), Streptococcus pneumoniae , and Enterococcus species, are examples of gram positive bacteria. An example of a Staphylococcus aureus species is methicillin-resistant Staphylococcus aureus (“MRSA”) and Enterococcus species include Vancomycin-resistant Enterococcus (“VRE”). Streptococcus pneumoniae is a causative pathogen of, for example, pneumonia and sepsis. SUMMARY Provided is a micro-device for separating and/or disrupting cells, which facilitates cell disruption. Provided is a method for disrupting cells using the micro-device for disrupting cells. Additional aspects will be set forth in part in the description which follows and, in part, will be apparent from the description, or may be learned by practice of the embodiments set forth herein. According to an embodiment of the invention, a micro-device for separating and/or disrupting cells includes a first chamber in which the cells are disrupted, a second chamber which is pressurized and depressurized, a flexible membrane which separates the first chamber and the second chamber and is vibrated by pressuring and depressurizing the second chamber, and a micro-unit confined in the first chamber, where the micro-unit disrupts the cells in the first chamber. In an embodiment, the micro-unit may include a plurality of solid supports. The solid support may be present as a separate material with respect to the first chamber wall, i.e., present in the inner space of the first chamber without permanently attached to the wall of the chamber. The solid support are not limited to microscale dimensions, and may have any dimensions so long as a cell can be separated and/or lysed by the operation of the cell lysis device. The solid support may be a microparticle. The term “microparticle” refers to a particle having the longest length in the cross-section of the solid support up to about 1000 μm. For example, the longest length may be from about 10 nanometers (nm) to about 1000 μm, specifically about 10 nm to about 700 μm, more specifically about 10 nm to about 500 μm, or about 10 nm to about 300 μm, or about 10 nm to about 100 μm, or about 10 nm to about 50 μm, or about 10 nm to about 10 μm in size, and in an embodiment, may be from about 50 nm to about 1000 μm, specifically about 50 nm to about 900 μm, more specifically about 50 nm to about 700 μm, or about 50 nm to about 500 μm, or about 50 nm to about 300 μm, or about 50 nm to about 100 μm, or about 50 nm to about 50 μm, or about 50 nm to about 10 μm in size, and in another embodiment, may be from about 100 nm to about 1000 μm, specifically about 100 nm to about 900 μm, more specifically about 100 nm to about 700 μm, or about 100 nm to about 500 μm, or about 100 nm to about 300 μm, or about 100 nm to about 100 μm, or about 100 nm to about 50 μm, or about 100 nm to about 10 μm, and in yet another embodiment, may be from about 1 μm to about 500 μm, specifically about 1 μm to about 300 μm, more specifically about 1 μm to about 100 μm, or about 1 μm to about 50 μm, or about 1 μm to about 30 μm, or about 1 μm to about 10 μm in size. The solid support may be at least one of spherical, planar, or multi-planar. The solid support may be microbeads. The solid supports in the first chamber may be at least one, for example two in number. In an embodiment, the number of solid supports in the first chamber may be 10 or greater, specifically 100 or greater, more specifically 1000 or greater, or 10,000 or greater, or 100,000 or greater, or 10 8 or greater. For example, the number of solid supports in the first chamber may be from about 1 to about 10 8 , specifically about 100 to about 10 6 . The first chamber may be manufactured such that the solid supports are placed in a space and blocked by the first chamber, so that the solid supports, for example beads may be unable to pass out of an inlet or an outlet. A density D of the solid supports in the first chamber may be greater than or equal to 1 gram per cubic centimeter (g/cm 3 ), specifically about 1 g/cm 3 to about 20 g/cm 3 , more specifically about 1 g/cm 3 to about 15 g/cm 3 , or about 1 g/cm 3 to about 10 g/cm 3 , or about 1 g/cm 3 to about 8 g/cm 3 , or about 1 g/cm 3 to about 6 g/cm 3 , or about 1 g/cm 3 to about 4 g/cm 3 , or about 3 g/cm 3 to about 20 g/cm 3 , or about 3 g/cm 3 to about 15 g/cm 3 , or about 3 g/cm 3 to about 10 g/cm 3 , or about 3 g/cm 3 to about 8 g/cm 3 , or about 3 g/cm 3 to about 6 g/cm 3 , or about 3 g/cm 3 to about 4 g/cm 3 . In an embodiment, the first chamber may include an inlet, through which the cells are introduced, and an outlet, through which contents of disrupted cells are released. In an embodiment, the second chamber may include a plurality of sub-chambers. In an embodiment, the solid support may be a plurality of microbeads and the plurality of microbeads may include at least one of a glass bead, a metallic bead and a metal oxide bead. In an embodiment, the micro-unit may further include a binding material which binds to one or more cells or viruses on its surface. The binding material may an organic layer disposed on a surface of the plurality of solid support for example, microbeads. In an embodiment, at least one of a wall of the first chamber and a surface of the solid supports such as microparticles may comprise a binding material which is able to bind to a cell or virus, and in an embodiment, may be coated with the binding material. The binding material may be a material which specifically binds to the cell or virus. The specific binding material may be at least one of an antibody for an antigen, a substrate or inhibitor for an enzyme, an enzyme for a substrate, a receptor for a ligand, and a ligand for a receptor. The binding material may be a material which non-specifically binds to the cell or virus. The nonspecific binding material may be a hydrophobic material having a water contact angle of about 70° to about 95°, and in an embodiment, may include a material with at least one amino group. The material with at least one amino group may include a polymer having at least two amino groups. The material with at least one amino group may include polyethyleneimine (“PEI”). The hydrophobic material with a water contact angle of about 70° to about 95° may include octadecyldimethyl(3-trimethoxysilyl propyl)ammonium (“OTC”) or tridecafluorotetrahydrooctyltrimethoxylsilane (“DFS”). The wall of the first chamber or the surface of the solid supports comprising the binding material may be of use in capturing or adsorbing the cell or virus to separate. The material with at least one amino group may an aminosilane moiety having at least one amino group selected from primary amino group, secondary amino group, tertiary amino group, and quaternary amino group. The tertiary amino group may be any tertiary amino group, except that a amide group or a nitrile group may not be used when used alone. The aminosilane moiety may be derived from an aminosilane molecule by depositing the surface of the nonplanar solid support. The aminosilane molecule may be an organic aminosilane molecule having a formula (X 1 )(X 2 )(X 3 )Si(Y), wherein X 1 , X 2 and X 3 are independently selected from the group consisting of hydrogen, alkoxy group (—OR) and halogen and at least one of X 1 , X 2 and X 3 is alkoxy group. In the alkoxy group (—OR), R may be a hydrocarbonyl group having one to 20 carbon atoms, for example, methyl, ethyl, propyl, iso-propyl etc. The halogen may be F, Cl, Br, I or At. Y may an organic moiety containing at least one amino group. The organic moiety may be aminoalkyl, or polyethyleneimine. In the aminoalkyl, the alkyl group may have one to 20 carbons atoms. The polyethyleneimine may have the formula —[CH 2 CH 2 NH] n —, wherein n may be 2-100. The alkoxy groups (—OR) hydrolyze in an aqueous environment, producing hydroxyl group, one or more of which may undergo condensation and elimination reactions with surface —OH groups found on the surface of the solid support as well as on neighboring organosilane molecules. The aminosilane molecule may 3-aminopropyltriethoxysilane (GAPS), polyethyleneiminetrimethoxysilane such as N-(2-aminoethyl)-3-aminopropyltriethoxysilane (EDA) and (3-trimethoxysilyl-propyl)diethylenetriamine (DETA). The aminosilane molecules may be coated on the surface of the solid support by a known method. For example, the coating may be conducted by using dip-coating method or chemical vapor deposition (CVD) method. The dip-coating method may be conducted according to following exemplary procedure: solid supports such as glass substrates may be cleaned using a general glass cleaning protocol: immersion in NaOH (2.5 M) solution for 24 h, sonication in H 2 O for 10 min, immersion in HCl (0.1 M) for 15 min, sonication in H 2 O for 10 min and immersion in methanol for 5 min prior to the silanization step. Silanization of the samples may be performed by dip coating in 1% aqueous solutions of silane for 15 min. Post-treatment steps may include shaking in methanol for 5 min, rinsing in H 2 O for 10 min and finally spin drying (2000 rpm) for 5 min. The coated slides may be baked at 110° C. for 15 min and stored in a vacuum desiccator prior to analysis. Further, the dip-coating method may be conducted according to the method described in the examples. The CVD method may be conducted according to following exemplary procedure: solid supports such as a silicon wafer or the glass substrates. Using a pulsed CVD technique (ThinSonic pulsed ultrasonic CVD; Sono-Tech Corp.), different aminosilane molecule film such as GAPS film thicknesses may be deposited on the solid supports such as a silicon wafer or glass substrates. The pulsed ultrasonic method may be a process in which the precursor is delivered to the ultrasonic nozzle through a series of automatically controlled solenoid valves. The experimental apparatus may be known in the art. The precursor may atomized at the tip of the nozzle and introduced into a low-pressure reaction chamber in a near-vapor-phase state (about 15 μm drop diameter). Nozzle power, volume of pulses and number of pulses may be varied to obtain the desired coating thicknesses. All of the solid supports such as a silicon wafer or the glass substrates may be first cleaned in NH 4 OH:H 2 O 2 :H 2 O (1:1:5) solution before the silane deposition step. The substrates may be heated (120° C.) in situ for 15 min prior to the pulses injections. The coated solid support may be dried under N2 without any post-washing steps. In this application, the term “water contact angle” refers to water contact angle measured by a Kruss Drop Shape Analysis System type DSA 10 Mk2. A droplet of 1.5 ul deionized water is automatically placed on the sample. The droplet was monitored every 0.2 seconds for a period of 10 seconds by a CCD-camera and analyzed by Drop Shape Analysis software (DSA version 1.7, Kruss). The complete profile of the droplet was fitted by the tangent method to a general conic section equation. The angles were determined both at the right and left side. An average value is calculated for each drop and a total of five drops per sample are measured. The average of the five drops is taken the contact angle. In an embodiment, a diameter of each of the plurality of spherical solid supports may be in a range of about 1 micrometer to about 500 micrometers. In an embodiment, a density of the plurality of solid supports such as microbeads may be in a range of about 1 gram per cubic centimeter to about 20 grams per cubic centimeter. In an embodiment, the solid supports such as beads may be rigid enough to be used as a cell lysis media. Solid supports, as used herein, refer to spherical bodies or non-spherical bodies, and can have any shape so long as the shape does not adversely affect the desirable features of the disclosed method or device. The solid supports may comprise a solid material at least a part thereof, and in an embodiment, may include a coating with a solid material. In an embodiment, the solid supports are a solid particle. In another embodiment, the beads are hollow. The beads may be magnetic or nonmagnetic. The beads may comprise at least one of a glass, metal, and metal oxide. In an embodiment, the beads comprise at least one of glass beads, metal beads, and metal oxide beads. The metal oxide may comprise at least one of ZrO 2 , SiO 2 , Al 2 O 3 , Fe 2 O 3 , and TiO 2 . In an embodiment, the solid supports comprise a combination of, for example, ZrO 2 and SiO 2 . The metal beads may comprise, for example, steel beads or stainless steel beads. Glass beads are specifically mentioned. In an embodiment, a size of at least one of the inlet and the outlet may be less than a size of the micro-unit. In an embodiment, a size of at least one of the inlet and the outlet may be greater than a size of the micro-unit, and the first chamber may further include at least one projection on an inner surface of the at least one of the inlet and the outlet, where the at least one projection makes an effective size of the at least one of the inlet and the outlet less than a size of the micro-unit. In an embodiment, a thickness of the membrane may be in a range of about 1 micrometer to about 5 millimeters. The membrane may be flexible and elastic and thus, may be extended to upward or downward by pressuring and depressurizing the second chamber. The flexible and elastic membrane may be placed between the first chamber and the second chamber and define at least a part of the wall of the first chamber and second chamber so as to be able to be vibrated by pressuring and depressurizing the second chamber. The membrane may comprise an elastic polymer membrane, such as a polydimethylsiloxane (“PDMS”) membrane. In an embodiment, the cells may include at least one of bacteria, virus and fungi. According to another embodiment of the invention, a method of disrupting cells includes: introducing cells into a first chamber of a micro-device for disrupting cells, where the micro-device includes the first chamber, a second chamber which is pressurized and depressurized, a flexible membrane which separates the first chamber and the second chamber and is vibrated by pressurizing and depressurizing the second chamber, and a micro-unit for disrupting cells confined in the first chamber; disrupting the cells by vibrating the flexible membrane between the first chamber and the second chamber, after the introducing the cells; and releasing resultant disrupted cells from the first chamber of the micro-device. In an embodiment, the introducing the cells may include passing a solution containing the cells through the first chamber. In an embodiment, the method may further include providing a cell lysis solution to the first chamber. In an embodiment, the disrupting of the cells may include vibrating the plurality of micro-units by vibrating the flexible membrane which separates the first chamber from the second chamber. The micro-units are as described herein above. The micro-units such as a solid support for disrupting cells are contained in the first chamber and may be present as a separate material with respect to the chamber wall, i.e., present in the inner space of the first chamber without permanently attached to the wall of the chamber. In an embodiment, the cell lysis solution may include at least one of NaOH, KOH, chaotrope, a surfactant, cell wall degrading enzyme and a biological buffer. In an embodiment, the disrupting the cells may include periodically or non-periodically adjusting the pressure of the interior of the second chamber such that the flexible membrane is vibrated. In an embodiment, the disrupting the cells may include vibrating the flexible membrane at a range of about 0.001 hertz to about 100 kilohertz. In an embodiment, the method may further include disrupting the cells by vibrating the flexible membrane between the first chamber and the second chamber, after the introduction of the cells. The membrane may be flexible and elastic and thus, may be extended to upward or downward by pressuring and depressurizing the second chamber. The flexible and elastic membrane may be placed between the first chamber and the second chamber and define at least a part of the wall of the first chamber and second chamber so as to be able to be vibrated by pressuring and depressurizing the second chamber. In an embodiment, the micro-unit may include a microbead. In an embodiment, the step of introducing the cells or viruses may comprise passing a solution containing the cells or viruses at a rate of 100 to 500 ul/min through the first chamber such that the cells or viruses in the solution is retained on the solid support and the first chamber has a volume of 1 ul to 100 ul and the solution has a pH of 3.0 to 6.0 and a salt concentration of 10 mM to 500 mM, wherein the solid support comprises a binding material which binds to a cell or a virus on its surface. The binding material may be a material which specifically binds to the cell or virus. The specific binding material may be at least one of an antibody for an antigen, a substrate or inhibitor for an enzyme, an enzyme for a substrate, a receptor for a ligand, and a ligand for a receptor. The binding material may be a material which non-specifically binds to the cell or virus. The nonspecific binding material may be a hydrophobic material having a water contact angle of about 70° to about 95°, and in an embodiment, may include a material with at least one amino group. The material with at least one amino group may include a polymer having at least two amino groups. The material with at least one amino group may include polyethyleneimine (“PEI”). The hydrophobic material with a water contact angle of about 70° to about 95° may include octadecyldimethyl(3-trimethoxysilyl propyl)ammonium (“OTC”) or tridecafluorotetrahydrooctyltrimethoxylsilane (“DFS”). The wall of the first chamber or the surface of the solid supports comprising the binding material may be of use in capturing or adsorbing the cell or virus to separate. The material with at least one amino group may an aminosilane moiety having at least one amino group selected from primary amino group, secondary amino group, tertiary amino group, and quaternary amino group. The tertiary amino group may be any tertiary amino group, except that a amide group or nitrile group may not be used when used alone. The aminosilane moiety may be derived from an aminosilane molecule by depositing the surface of the nonplanar solid support. The aminosilane molecule may be an organic aminosilane molecule having a formula (X 1 )(X 2 )(X 3 )Si(Y), wherein X 1 , X 2 and X 3 are independently selected from the group consisting of hydrogen, alkoxy group (—OR) and halogen and at least one of X 1 , X 2 and X 3 is alkoxy group. In the alkoxy group (—OR), R may be a hydrocarbonyl group having one to 20 carbon atoms, for example, methyl, ethyl, propyl, iso-propyl etc. The halogen may be F, Cl, Br, I or At. Y may an organic moiety containing at least one amino group. The organic moiety may be aminoalkyl, or polyethyleneimine. In the aminoalkyl, the alkyl group may have one to 20 carbons atoms. The polyethyleneimine may have the formula —[CH 2 CH 2 NH] n —, wherein n may be 2-100. The alkoxy groups (—OR) hydrolyze in an aqueous environment, producing hydroxyl group, one or more of which may undergo condensation and elimination reactions with surface —OH groups found on the surface of the solid support as well as on neighboring organosilane molecules. The aminosilane molecule may 3-aminopropyltriethoxysilane (GAPS), polyethyleneiminetrimethoxysilane such as N-(2-aminoethyl)-3-aminopropyltriethoxysilane (EDA) and (3-trimethoxysilyl-propyl)diethylenetriamine (DETA). The aminosilane molecules may be coated on the surface of the solid support by a known method. For example, the coating may be conducted by using dip-coating method or chemical vapor deposition (CVD) method. The dip-coating method may be conducted according to following exemplary procedure: solid supports such as glass substrates may be cleaned using a general glass cleaning protocol: immersion in NaOH (2.5 M) solution for 24 hours, sonication in H2O for 10 min, immersion in HCl (0.1 M) for 15 min, sonication in H2O for 10 min and immersion in methanol for 5 min prior to the silanization step. Silanization of the samples may be performed by dip coating in 1% aqueous solutions of silane for 15 min. Post-treatment steps may include shaking in methanol for 5 min, rinsing in H2O for 10 min and finally spin drying (2000 rpm) for 5 min. The coated slides may be baked at 110° C. for 15 min and stored in a vacuum desiccator prior to analysis. The CVD method may be conducted according to following exemplary procedure: solid supports such as a silicon wafer or the glass substrates. Using a pulsed CVD technique (ThinSonic pulsed ultrasonic CVD; Sono-Tech Corp.), different aminosilane molecule film such as GAPS film thicknesses may be deposited on the solid supports such as a silicon wafer or glass substrates. The pulsed ultrasonic method may be a process in which the precursor is delivered to the ultrasonic nozzle through a series of automatically controlled solenoid valves. The experimental apparatus may be known in the art. The precursor may atomized at the tip of the nozzle and introduced into a low-pressure reaction chamber in a near-vapor-phase state (about 15 μm drop diameter). Nozzle power, volume of pulses and number of pulses may be varied to obtain the desired coating thicknesses. All of the solid supports such as a silicon wafer or the glass substrates may be first cleaned in NH4OH:H2O2:H2O (1:1:5) solution before the silane deposition step. The substrates may be heated (120° C.) in situ for 15 min prior to the pulses injections. The coated solid support may be dried under N2 without any post-washing steps. In this application, the term “water contact angle” refers to water contact angle measured by a Kruss Drop Shape Analysis System type DSA 10 Mk2. A droplet of 1.5 mu.l deionized water is automatically placed on the sample. The droplet was monitored every 0.2 seconds for a period of 10 seconds by a CCD-camera and analyzed by Drop Shape Analysis software (DSA version 1.7, Kruss). The complete profile of the droplet was fitted by the tangent method to a general conic section equation. The angles were determined both at the right and left side. An average value is calculated for each drop and a total of five drops per sample are measured. The average of the five drops is taken the contact angle. In an embodiment of the invention, the method may further comprise washing the solid support with wash buffer by flowing the wash buffer at a rate of 100 to 500 ul/min through the first chamber to remove impurities form the solid support, wherein the wash buffer has a pH of 3.0 to 9.0, for example 3.0 to 6.0, 3.0 to 7.0 or 3.0 to 8.0, and a salt concentration of 10 mM to 500 mM. The wash buffer may be acetate buffer or phosphate buffer or TRIS buffer. The solution comprising cells or viruses comprises both animal cells and bacterial cells, and the impurities are animal cells, thereby after the step of washing, the animal cells are selectively removed from the solid support and retaining bacterial cells to the solid support. In an embodiment of the invention, prior to the step of introducing, the method may further include diluting a sample comprising cells or viruses with a phosphate buffer or an acetate buffer having a pH of 3.0 to 6.0 and salt concentration of 10 mM to 500 mM to obtain a solution containing cells or viruses having a pH of 3.0 to 6.0 and salt concentration of 10 mM to 500 mM. In an embodiment, the method further includes performing polymerase chain reaction (PCR) using the lysed product obtained from the disrupting step as a template without any further purification step. The sample containing the cells or viruses may be blood or a sample containing the blood derived material such red blood cell, white blood cells, platelets etc. Thus, PCR amplification may be achieved with rapidity and convenience since no further step for removing the PCR inhibitor is required. The lysis solution may be NaOH for example, 0.0001 N to 0.1 N NaOH. According to another embodiment of the invention, a micro-device includes a membrane, a first chamber which is partly delimited by the membrane, a second in the first chamber, where the plurality of microparticles contacts the membrane, and a port in fluid communication with the second chamber, where the membrane is vibrated by pressurizing and depressurizing the second chamber. In an embodiment, a density of the plurality of microparticles may be in a range of about 1 gram per cubic centimeter to about 20 grams per cubic centimeter. BRIEF DESCRIPTION OF THE DRAWINGS These and/or other aspects will become apparent and more readily appreciated from the following description of the embodiments, taken in conjunction with the accompanying drawings in which: FIG. 1 is a cross sectional view of an embodiment of a micro-device for disrupting cells in a molecular diagnostic apparatus according to the invention; FIG. 2 is a cross section view of an embodiment of a bead including an organic layer thereon; FIGS. 3 to 6 are cross sectional views of alternative embodiments of the micro-device for disrupting cells in a molecular diagnostic apparatus according to the invention; FIG. 7 is a graph showing threshold cycle (number) versus actuation time (minute) representing results of a cell disruption test conducted according to a method of disrupting cells. FIGS. 8 and 9 illustrate a cell lysis device according an embodiment of the cell lysis device. FIG. 10 is a graph of threshold cycle (Ct) and change in Ct (ΔCt) versus liquid volume fraction (f L ) showing cell lysis efficiency with respect to changes in liquid volume fraction (f L ). DETAILED DESCRIPTION The invention now will be described more fully hereinafter with reference to the accompanying drawings, in which various embodiments are shown. This invention may, however, be embodied in many different forms, and should not be construed as limited to the embodiments set forth herein. Rather, these embodiments are provided so that this disclosure will be thorough and complete, and will fully convey the scope of the invention to those skilled in the art. Like reference numerals refer to like elements throughout. It will be understood that when an element is referred to as being “on” another element, it can be directly on the other element or intervening elements may be present therebetween. In contrast, when an element is referred to as being “directly on” another element, there are no intervening elements present. As used herein, the term “and/or” includes any and all combinations of one or more of the associated listed items. It will be understood that, although the terms first, second, third etc. may be used herein to describe various elements, components, regions, layers and/or sections, these elements, components, regions, layers and/or sections should not be limited by these terms. These terms are only used to distinguish one element, component, region, layer or section from another element, component, region, layer or section. Thus, a first element, component, region, layer or section discussed below could be termed a second element, component, region, layer or section without departing from the teachings of the present invention. The terminology used herein is for the purpose of describing particular embodiments only and is not intended to be limiting. As used herein, the singular forms “a,” “an” and “the” are intended to include the plural forms as well, unless the context clearly indicates otherwise. It will be further understood that the terms “comprises” and/or “comprising,” or “includes” and/or “including” when used in this specification, specify the presence of stated features, regions, integers, steps, operations, elements, and/or components, but do not preclude the presence or addition of one or more other features, regions, integers, steps, operations, elements, components, and/or groups thereof. Furthermore, relative terms, such as “lower” or “bottom” and “upper” or “top,” may be used herein to describe one element's relationship to another element as illustrated in the Figures. It will be understood that relative terms are intended to encompass different orientations of the device in addition to the orientation depicted in the Figures. For example, if the device in one of the figures is turned over, elements described as being on the “lower” side of other elements would then be oriented on “upper” sides of the other elements. The exemplary term “lower,” can therefore, encompasses both an orientation of “lower” and “upper,” depending on the particular orientation of the figure. Similarly, if the device in one of the figures is turned over, elements described as “below” or “beneath” other elements would then be oriented “above” the other elements. The exemplary terms “below” or “beneath” can, therefore, encompass both an orientation of above and below. Unless otherwise defined, all terms (including technical and scientific terms) used herein have the same meaning as commonly understood by one of ordinary skill in the art to which this invention belongs. It will be further understood that terms, such as those defined in commonly used dictionaries, should be interpreted as having a meaning that is consistent with their meaning in the context of the relevant art and the present disclosure, and will not be interpreted in an idealized or overly formal sense unless expressly so defined herein. Exemplary embodiments are described herein with reference to cross section illustrations that are schematic illustrations of idealized embodiments. As such, variations from the shapes of the illustrations as a result, for example, of manufacturing techniques and/or tolerances, are to be expected. Thus, embodiments described herein should not be construed as limited to the particular shapes of regions as illustrated herein but are to include deviations in shapes that result, for example, from manufacturing. For example, a region illustrated or described as flat may, typically, have rough and/or nonlinear features. Moreover, sharp angles that are illustrated may be rounded. Thus, the regions illustrated in the figures are schematic in nature and their shapes are not intended to illustrate the precise shape of a region and are not intended to limit the scope of the present claims. All methods described herein can be performed in a suitable order unless otherwise indicated herein or otherwise clearly contradicted by context. The use of any and all examples, or exemplary language (e.g., “such as”), is intended merely to better illustrate the invention and does not pose a limitation on the scope of the invention unless otherwise claimed. No language in the specification should be construed as indicating any non-claimed element as essential to the practice of the invention as used herein. Hereinafter, embodiments of the invention will be described in further detail with reference to the accompanying drawings. FIG. 1 is a cross sectional view of an embodiment of a micro-device for disrupting cells in a molecular diagnostic apparatus, and FIG. 2 is a cross section view of an embodiment of a bead including an organic layer thereon. Referring to FIG. 1 , the micro-device 20 includes an upper plate 21 , a lower plate 23 and a membrane 26 . The membrane 26 is disposed between the upper plate 21 and the lower plate 23 . A space in a portion of the upper plate 21 defines a first chamber 22 , and a space in a portion of the lower plate 23 defines a second chamber 24 . The first chamber 22 and the second chamber 24 are separated by the membrane 26 , disposed between the first chamber 22 and the second chamber 24 . In such an embodiment, the first chamber 22 is delimited by the upper plate 21 and the membrane 26 , and the second chamber 24 is delimited by the lower plate 23 and the membrane 26 . In some embodiments, each of the first chamber and the second chamber may have a volume in a range from about 1 microliter (ul) to about 10,000 microliters (ul). The membrane 26 may be flexible. In an embodiment, the membrane 26 may be a polymer membrane, such as polydimethylsiloxane (“PDMS”) membrane, for example. The thickness of the membrane 26 may be, for example, in a range from about 1 micrometer (um) to about 5 millimeters (mm). The membrane 26 may be liquid and gas-permeable, but it may also be partly permeable. The first chamber 22 may include a plurality of micro-units, e.g., a plurality of particles. The micro-units may be magnetic or non-magnetic micro-units. In an embodiment, the particles may be a plurality of beads 28 , as shown in FIG. 1 . Since the first chamber 22 is delimited by the membrane 26 , the beads 28 may contact with the membrane 26 . In an embodiment, the beads 28 may be microbeads. In the exemplary embodiment shown in FIG. 1 , the micro-units include the plurality of beads 28 , but the invention is not limited thereto. In an alternative embodiment, other units may be included in the first chamber 22 instead of the plurality of beads 28 . In an embodiment, the diameter of each bead 28 may be, for example, in the range of about 1 um to about 500 um. The density of beads 28 in the first chamber 22 may be greater than about 1 gram per cubic centimeter (g/cm 3 ), for example, in the range of about 1 g/cm 3 to about 20 g/cm 3 . In an embodiment, at least one bead may be included in the liquid medium, for example, there may be 10, 100, 1000, 10,000, 10 5 , 10 6 , 10 7 , 10 8 , or 10 9 or more beads per 1 microliter (ul) of the liquid medium. In some embodiments, 1 to 10 8 beads, for example, 100 to 10 6 beads, may be included per 1 ul of the liquid medium in the first chamber 22 . The micro-units may have a sphere shape, a plate shape, or may have a shape including a plurality of sides. The micro-units may be magnetic or non-magnetic beads. In an embodiment, the beads 28 may be glass beads. In an alternative embodiment, the beads 28 may be metal oxide beads or metallic beads. The metal oxide may be one of ZrO 2 , SiO 2 , Al 2 O 3 , Fe 2 O 3 , TiO 2 and a mixture thereof. In one embodiment, the mixture may be, for example, a mixture including ZrO 2 and SiO 2 . The metal beads may be, for example, formed of steel or stainless steel. In an embodiment, the beads 28 have a composition of glass or metal oxide, and surface modifications for cell capture or absorption are thereby effectively implemented. In an embodiment, surfaces of the beads 28 may be capable of binding to a cell or may be modified to be suitable for cell capture by binding to a cell. The surfaces of the beads 28 may be hydrophobic, for example, have a water contact angle in a range of 70° to 95°. The hydrophobicity may be rendered, for example, by coating the surfaces of the beads 28 with compounds such as octadecyldimethyl(3-trimethoxysilylpropyl)ammonium (“OTC”) or tridecafluoro tetrahydrooctyltrimethoxysilane (“DFS”). The surfaces of the beads 28 may also be coated with polyethyleneimine trimethoxy silane (“PEIM”). The surfaces of the beads 28 may have cell-binding substances immobilized thereon, such as an antibody which binds to an antigen on a cell surface, a receptor which binds to a ligand on a cell surface, or a ligand which binds to a receptor on a cell surface, for example. In an embodiment, the surfaces of the beads 28 include substances capable of specifically binding to specific cells such that the beads 28 specifically separate specific cells. In an alternative embodiment, the surfaces of the beads 28 include substances capable of non-specifically binding to cells such that the beads 28 non-specifically separate all kinds of cells. The cells may include bacteria, a virus or fungi. The cells may be contained in an appropriate liquid medium. The liquid medium may be, for example, a medium for cell culture, a buffer, for example, phosphate buffered saline (“PBS”), physiological saline, body fluids or water. The liquid medium may also include a cell lysis solution. The cell lysis solution may be additionally and separately supplied to the chamber after the liquid medium containing the cells is supplied, or may be pre-mixed and then supplied to the chamber. The cell lysis solution may include a non-specific cell lysis agent and/or a specific cell lysis agent. The non-specific cell lysis agent may include at least one of a surfactant, NaOH and a chaotropic salt. The specific cell lysis agent may include, cell wall degrading enzyme for example, lysozyme, lysostaphin or penicillin and beta-lactam antibiotics. The cells bound to the micro-units may be washed, after the cells are introduced into the first chamber 22 . The washing may be carried out by passing a wash solution through the first chamber 22 with or without the liquid medium. The wash solution may wash off unbound substances, leaving the cells bound to the micro-units. The wash solutions may be, for example, water, a buffer (such as a PBS buffer) or physiological saline. In one embodiment, the cells may be disrupted while the cells are introduced into the first chamber 22 . In an alternative embodiment, the cells may be disrupted after the cells are introduced into the first chamber 22 . The micro-device 20 includes an inlet 30 and an outlet 32 . The sizes of the inlet 30 and outlet 32 may be smaller than a size of the micro-units, e.g., a diameter of each bead 28 . In one exemplary embodiment, the inlet 30 and the outlet 32 may have a circular cross-sectional shape, and the diameters of the inlet 30 and the outlet 32 may be less than the diameter of the bead 28 . A solution containing cells to be disrupted is introduced into the first chamber 22 through the inlet 30 . Resultant disrupted cells, including a nucleic acid, for example, obtained by the disruption of cell membrane and/or walls are released through the outlet 32 . The inlet 30 may be formed through one wall of the upper plate 21 to connect to one side of the first chamber 22 . The outlet 32 may be formed through another wall of the upper plate 21 to connect to the other side of the first chamber 22 . In an embodiment, at least one of the inlet 30 and outlet 32 may be operatively connected to a unit for providing power (not shown). The unit for providing power can provide a power to move a fluid through at least one of the inlet 30 and outlet 32 . The unit for providing power may include a unit causing the motion of fluid, for example, a unit that provides a positive pressure or a negative pressure to the first chamber 22 , including a pump. The pump may be a micropump, which may be applied to a microfluidic device. The micropump may be a mechanical or a non-mechanical device. The mechanical micropump may include an actuator and moving parts which are membranes or flaps. The motion of fluid may be generated using a piezoelectric, electrostatic, thermo-pneumatic, pneumatic or magnetic effect. A non-mechanical device may function as the unit for providing the power by generating an electro-hydrodynamic force, or an electro-osmotic or ultrasonic flow. The inlet 30 and the outlet 32 may be in fluid-communication with the first chamber 22 , for example, through a microchannel (not shown). The microchannel may have a width in a range of about 1 um to about 10,000 um, for example about 1 um to about 5,000 um. The first chamber 22 including the micro-units may be in fluid-communication with at least one of a storage unit (not shown) that stores a cell lysis solution and a storage unit (not shown) that stores cell wash solution. The storage units may be connected to the chamber through the inlet 30 . The cells may be introduced by applying a positive pressure to an inlet 30 of the first chamber 22 or applying a negative pressure to the outlet 32 of the first chamber 22 . In an embodiment, the negative pressure or positive pressure may be applied by a pump (not shown). The pump may be at least one of a peristaltic pump and a pneumatic pump. In an alternative embodiment, the cells may be introduced through direct infusion by a user. In one embodiment, for example, the cells may be infused by the user performing pipetting. The amount and rate of introduction may depend on the cells to be disrupted, the purpose of cell disruption, and the post-cell disruption process, for example, and one of ordinary skill in the art may appropriately adjust them. The application of the pressure may be carried out in a state in which both the inlet 30 and the outlet 32 are closed. That is, the cells may be disrupted in a closed chamber containing the cells. The application of the pressure may be also carried out in a state in which at least one of the inlet 30 and the outlet 32 is open. That is, the cells may be disrupted under conditions that at least a portion of the liquid medium containing the cells in the chamber is flowing. The second chamber 24 may operate as a pneumatic chamber including a space into which a fluid, such as air, for pressing periodically or non-periodically on the membrane 26 is introduced. High-pressurized fluid is introduced into the second chamber 24 , and thus the membrane 26 is pressed. When the membrane 26 is pressed, the membrane 26 protrudes toward the first chamber 22 , and the spatial volume of the first chamber 22 is thereby reduced. When the membrane 26 is de-pressed, the membrane 26 shrinks down toward the second chamber 24 . The second chamber 24 has a port 34 which is an inflow passage of the fluid that pressurizes the interior of the second chamber 24 and simultaneously is an outflow passage of the fluid. In an embodiment, the fluid may be periodically or non-periodically introduced into/the second chamber 24 through the port 34 or released from the second chamber 24 through the port 34 , and thus the membrane 26 may be periodically or non-periodically vibrated. The vibration of the membrane 26 leads to a periodic or non-periodic pressure to the beads 28 in the first chamber 22 through direct contact with the beads 28 or through the solution contained in the first chamber 22 . In such an embodiment, motion of the beads 28 is induced, and the beads 28 collide with each other or collide with an inner wall of the first chamber 22 . Due to the motion of the beads 28 , the cells introduced into the interior of the first chamber 22 are disrupted by being sheared or grinded, that is the cells may be disrupted by a shearing force or an impact force applied to the cell or by heat, which are induced by the motion of the micro-units. The pressurizing or depressurizing of the interior of the second chamber 24 through the introduction of the fluid into the second chamber 24 (applying a positive pressure) or release of the fluid from the second chamber 24 (applying a negative pressure) is effectively controlled with a vibration frequency in a range of from about 0.001 hertz (Hz) to about 100 kilohertz (kHz). Referring to FIG. 2 , surfaces of the beads 28 may have an organic layer 28 A which allows for specific or non specific cell capture. In one embodiment, for example, an antibody or an aptamer, may be coated on the surfaces of the beads 28 to selectively capture a specific cell. In an embodiment, a nonspecific cell may be captured by a hydrophobic or electrostatic force. The organic layer 28 A may be formed by modifying the surfaces of beads 28 in various ways using organosilane. In one embodiment, a portion having an affinity with a specific or nonspecific cell is on a surface of the beads 28 , and the cells introduced into the first chamber 22 are captured by the portion on the surface of the beads 28 . In such an embodiment, the vibration of the membrane 26 causes the motion of the beads 28 , which makes beads collide with each other or with the inner surface of the wall of the first chamber 22 , and the cells on the surface of the beads 28 may be disrupted by the collisions. In an alternative embodiment, after supplying a solution containing the cells to be disrupted, a substance for increasing the cell disruption effect may be supplied in the first chamber 22 , and then the cell disruption may be carried out. In an embodiment, the cell lysis solution may include the substance for increasing the cell disruption effect. The cell lysis solution may include, for example, NaOH, KOH, a chaotropic solution or a surfactant. In an embedment, the cell lysis solution may include a biological buffer such as Tris, phosphate, citrate, acetate and carbonate, for example, which does not increase the cell disruption effect. The cell lysis solution may be used in a concentration not affecting the post-cell disruption processes such as polymerase chain reaction (“PCR”), and thus PCR may be conducted after cell disruption without a further purification process. In such an embodiment, the substance is used in a concentration affecting PCR, and then a purification process is carried out. The cell lysis solution may be supplied after the cell disruption process, thereby facilitating the release of a nucleic acid. In an alternative embodiment, a solution containing the cells may be supplied to the first chamber 22 and then the cell disruption may be carried out without additional supply of the cell lysis solution into the chamber 22 . FIGS. 3 to 6 are cross sectional views of alternative embodiments of the micro-device for disrupting cells in a molecular diagnostic apparatus according to the invention. FIG. 3 shows an alternative embodiment of the micro-device 20 of FIG. 1 . The micro-device in FIG. 3 is substantially the same as the micro-device shown in FIG. 1 except for the inlet 30 and the outlet 32 . The same or like elements shown in FIG. 3 have been labeled with the same reference characters as used above to describe the embodiment of the micro-device shown in FIG. 1 , and any repetitive detailed description thereof will hereinafter be omitted or simplified. Referring to FIG. 3 , the sizes of the inlet 30 and the outlet 32 may be larger than a size of the bead 28 . In one exemplary embodiment, the inlet 30 and the outlet 32 may have a circular cross-sectional shape, and the diameters of the inlet 30 and the outlet 32 may be greater than the diameter of the bead 28 . A plurality of first projections 40 is disposed on an inner side of the inlet 30 . The first projections 40 may be evenly distributed throughout the inner side of the inlet 30 . The first projections 40 may be disposed in opposite direction with each other. Due to the first projections 40 , the substantial size or effective size of the inlet 30 , e.g., a size of a cross-sectional shape of the inlet 30 becomes smaller than the size of the beads 28 . Similarly, a plurality of second projections 42 is disposed the inner side of the outlet 32 . The distribution of the second projections 42 may be substantially the same as the distribution of the first projections 40 . Due to the second projections 42 , the substantial size or effective size of the outlet 32 , e.g., a size of a cross-sectional shape of the outlet 32 , may be smaller than the size of each bead 28 . In an embodiment, the shapes of the first and second projections 40 and 42 may be substantially identical to each other. In an alternative embodiment, the shapes of the first and second projections 40 and 42 may be substantially different from each other. In an embodiment, the first and second projections 40 and 42 may be formed by embossing the inner sides of the inlet 30 and the outlet 32 . FIG. 4 shows another alternative embodiment of the micro-device 20 of FIG. 1 . The micro-device in FIG. 4 is substantially the same as the micro-device shown in FIG. 1 except for the inlet 30 and the outlet 32 . The same or like elements shown in FIG. 4 have been labeled with the same reference characters as used above to describe the embodiment of the micro-device shown in FIG. 1 , and any repetitive detailed description thereof will hereinafter be omitted or simplified. Referring to FIG. 4 , the inlet 30 may be substantially the same as the inlet 30 of the embodiment shown in FIG. 3 . In an embodiment, a filter 44 may be disposed in the outlet 32 . The filter 44 may be a porous material which allows contents of the disrupted cells to pass. The size of the inlet 30 may be smaller than the size of the micro-unit, e.g., the size of each bead 28 , as in FIG. 1 . FIG. 5 shows another alternative embodiment of the micro-device 20 of FIG. 1 . The micro-device in FIG. 5 is substantially the same as the micro-device shown in FIG. 1 except for the second chamber 24 . The same or like elements shown in FIG. 5 have been labeled with the same reference characters as used above to describe the embodiment of the micro-device shown in FIG. 1 , and any repetitive detailed description thereof will hereinafter be omitted or simplified Referring to FIG. 5 , the second chamber 24 includes two chambers, e.g., a first sub-chamber 24 A and a second sub-chamber 24 B. The first sub-chamber 24 A and the second sub-chamber 24 B are separated by a partition wall 48 . The role of the third and second sub-chambers 24 A and 24 B may be the same as that of the second chamber 24 . The first sub-chamber 24 A includes a first port 34 A, and the second sub-chamber 24 B includes a second port 34 B. The structure of the first port and second port 34 A, 34 B may be substantially the same as the structure of the port 34 in the second chamber 24 of the embodiment shown in FIG. 1 . The structural of the inlet 30 and the outlet 32 may be substantially the same as to the inlet and the outlet 32 of the embodiments shown in FIG. 3 or FIG. 4 . The pressure may be applied to the first and second ports 34 A and 34 B, simultaneously or sequentially such that the membrane 26 of each sub-chamber vibrates simultaneously or sequentially. In an embodiment, the pressure may be applied with a same phase of pressure or different phases of pressure to the first and second ports 34 A and 34 B such that the first and second sub-chambers 24 A and 24 B vibrate in the same phase or different phases. In one embodiment, for example, the positive pressure is applied to the first port 34 A and the negative pressure is applied to the second port 34 B, and the membrane 26 of the first sub-chamber 24 A and the membrane 26 of the second sub-chamber 24 B thereby vibrate in opposite phases. FIG. 6 shows another alternative embodiment of the micro-device 20 . The micro-device in FIG. 6 is substantially the same as the micro-device shown in FIG. 1 except for the second chamber 24 . The same or like elements shown in FIG. 6 have been labeled with the same reference characters as used above to describe the embodiment of the micro-device shown in FIG. 1 , and any repetitive detailed description thereof will hereinafter be omitted or simplified Referring to FIG. 6 , the second chamber 24 includes three sub-chambers, i.e., first to third sub-chambers 24 A, 24 B and 24 C. The first to third sub-chambers 24 A, 24 B and 24 C may perform a function substantially the same as the function that the second chamber 24 in FIG. 1 performs. The first sub-chamber and the third sub-chamber 24 A and 24 C are separated by a first partition wall 48 A. The second sub-chamber and the third sub-chamber 24 B and 24 C are separated by a second partition wall 48 B. The first to third sub-chambers 24 A, 24 B and 24 C include first to third ports 34 A, 34 B and 34 C, respectively. The structure and function of the first to third ports 34 A 34 B, and 34 C may be substantially the same as the structure and function of the port 34 in FIG. 1 . The second chamber 24 of the embodiment in FIG. 6 includes three sub-chambers, but the invention is not limited thereto. In an alternative embodiment, the second chamber 24 may include more than three sub-chambers. The pressure may be applied to the first to third ports 34 A, 34 B and 34 C simultaneously or sequentially, allowing the membrane 26 in each chamber to vibrate simultaneously or sequentially. Also, the pressure may be applied at a same or different phase of pressure to the first to third ports 34 A, 34 B and 34 C, enabling the membrane 26 of each chamber to vibrate in the same or different phase. In one embodiment, for example, the positive pressure may be applied to the first port 34 A and the third port 34 C, and the negative pressure may be applied to the second port 34 B, such that the membrane 26 in the first sub-chamber and the third sub-chamber 24 A and 24 C and the membrane 26 in the second sub-chamber 24 B vibrate in opposite phases to one another. In the embodiments shown in FIGS. 1 and 3 to 6 , the first chamber 22 is disposed above the second chamber 24 , e.g., the first chamber 22 is an upper chamber and the second chamber 24 is a lower chamber, the invention is not limited thereto. In an alternative embodiment, the second chamber 24 may be disposed above the first chamber 22 , e.g., the second chamber 24 may be the upper chamber and the first chamber 22 may be the lower chamber. The operations of the micro-device where the second chamber 24 is disposed above the first chamber 22 is substantially the same as the embodiments shown in FIGS. 1 and 3 to 6 where the first chamber 22 is disposed above the second chamber 24 . Hereinafter, the invention will be described with reference to the following examples. It should be understood that the exemplary embodiments described herein should be considered in a descriptive sense only and not for purposes of limitation. EXAMPLE 1. Manufacture of a Micro-Device A micro-device having three layers (e.g., glass layer-PDMS layer-glass layer) was manufactured. A channel and a chamber were formed on the glass wafer through conventional photographic and etching processes and a wet etching process. After cleaning a 6 inch glass wafer (borosilicate glass, 700 um thickness) in a Piranha solution, e.g., a mixture of sulfuric acid (H 2 SO 4 ) and hydrogen peroxide (H 2 O 2 ), an amorphous polysilicon layer was vapor-deposited on the cleaned glass wafer to a thickness of 2000 um. Then, a patterning process was implemented, in which a part of the vapor-deposited polysilicon layer is exposed using a photoresist film. Then, the exposed part of the polysilicon layer was removed by dry etching. Thereafter, the photoresist film was stripped, and the exposed glass wafer was wet etched with a hydrofluoric acid solution (HF, 49%) to form a channel having a depth of 100 um and a width of 100 um. In the etching process to form the channel, a weir (projection) in a length of about 20 um was formed for isolating a bead. Next, the polysilicon layer was removed, and a dry film resist was coated and patterned. Then, a chamber (ca. 15.5 uL) for enclosing beads and holes were formed using a sand-blasting method. Thereafter, the glass wafer was diced into chip-shaped pieces, and cleaned with plasma. Then, a fluidic chip including the above-manufactured chamber in which beads are to be enclosed and a pneumatic chip including a chamber not enclosing the beads but functioning as a pneumatic pump were coupled to each other by way of a PDMS membrane (250 um thick) as an intermediate layer. FIG. 3 . shows the micro-device manufactured in the present example. About 15 mg to about 16 mg (about 2×10 5 in number) surface-modified glass beads were put into the produced bead chamber, the chamber was sealed using a tape, and the tape was covered with a plastic substrate to prevent the bending thereof upon operation. The surface-modified glass beads may be directly put into the bead chamber. Pressurizing and depressurizing the pneumatic chamber for vibrating the PDMS membrane and the motion of solution were controlled by a solenoid valve, an electro-regulator and LabVIEW software. 2. Modification of a Glass Bead After cleaning a glass bead (diameter: about 30 um to about 50 um, Polysciences, Inc.) in a Piranha solution, e.g., a mixture of sulfuric acid (H 2 SO 4 ) and hydrogen peroxide (H 2 O 2 ), the glass bead was sufficiently washed with distilled water and filtered under vacuum until dry. The glass bead may be spherical. Next, the bead was put into ethanol having 5 volume percent (% v/v) trimethoxysilylpropyl-modified (polyethyleneimine) (MW 1,500-1,800 Da, Cat #: SSP-060) (Gelest, Inc.) and a reaction was allowed to take place while mixing. After about 2 hours, the glass bead was sufficiently washed with an ethanol solution and filtered under vacuum until dry. Then it was sintered in an oven at 110° C. for 40 minutes. As a result, a glass bead having a surface coated with polyethyleneimine (“PEIM”) was obtained. Since the PEIM is capable of non-specifically binding to a cell, the glass bead may be employed for non-specifically separating a cell. 3. Experiments for Extracting a Nucleic Acid Gram-positive bacteria of S. aureus (“SA”) were diluted in a sodium acetate buffer (50 mM, pH 4) to a concentration of 10 6 colony-forming units per milliliter (CFU/ml). Then, 1 mL of the SA solution was introduced into a bead chamber by flowing through a channel, which is an inlet, to the chamber, at a flow rate of 200 ul/min for 5 minutes. The process was conducted with both the inlet and outlet of the chamber open. Then, a Tris buffer (Tris, pH 8, 10 mM) for washing the bead chamber was flowed through the channel and the chamber, at a flow rate of 200 microliters per minute (ul/min), and subsequently air was injected for drying the beads. Lastly, after injecting 0.02 N NaOH solution (6 uL) into the bead chamber at a very low rate, e.g., 30 kilopascal (kPa), the inlet and outlet were closed, and cells were disrupted by vibrating the PDMS membrane at a frequency of 10 Hz with +80 kPa and −80 kPa, by adjusting the Solenoid valve by means of the LabVIEW program. After the cell disruption process, the inlet and outlet were opened, 100 kPa of pressure was applied, 14 ul of NaOH solution was additionally injected by flowing through the chamber, and then the cell disruption product was recovered through the outlet. Thus, a cell disruption product (i.e, lysate) containing nucleic acid in a total of 20 ul of NaOH solution was obtained. For a positive control sample, a benchtop bead beating process was conducted. As the positive control experiment, 1 ml of SA dilute solution, which was prepared for the subject experiment, was centrifuged at 13200 revolutions per minute (rpm) for 20 minutes to precipitate SA bacteria, and then the supernatant was removed. A volume of 20 μL of 0.02N NaOH, the same cell lysis solution used in the device, and glass beads were put into a container including the precipitated bacteria and mixing was carried out with a vortexer (GENIE 2, Fisher) at maximum speed (Max. 3200 rpm) for 5 minutes, to obtain a cell disruption product. For a negative control, the supernatant was removed after centrifugation, deionized (“DI”) water alone was put into the container including the precipitated bacteria, and mixing was carried out with the above vortexer at the maximum speed for 5 minutes. 4. Real-Time Polymerase Chain Reaction (“PCR”) The amount of DNA extracted from the cell disruption achieved using the embodiment of the device disclosed herein for actuation times up to 20 min was compared with the amount of DNA extracted from the positive control sample by means of PCR (TMC PCR machine, Samsung). An SA 442 region present in S. aureus was tested using the following primers (Applied Biosystems, US), probes and composition of the PCR. TABLE 1 Sequence Sa442 forward 5′-GTT GCA TCG GAA ACA TTG TGT primer T-3′(SEQ. ID. No. 1) Sa442 reverse 5′-ATG ACC AGC TTC GGT ACT ACT primer AAA GAT-3′(SEQ. ID. No. 2) Sa442 Taqman 5′-TGT ATG TAA AAG CCG TCT probe TG-3′ (SEQ. ID. No. 3) Component volume(μl) Final concentration 10× Z-Taq buffer 0.2  1× 25 mM dNTP 0.16    2 mM Z-Taq polymerase 0.02 0.5 U 50 mM forward 0.03    1 μM primer 50 mM reverse 0.03    1 μM primer 20 mM probe 0.03  0.4 μM Water 0.5 — Disrupted cell 1 — solution 5. Experimental Results The implementation of the cell disruption device and method for five minutes resulted in a similar threshold cycle (Ct) to that obtained in the positive control experiment, as shown in FIG. 7 . Based on the experimental results shown in FIG. 7 , it may be concluded that an embodiment of the method of disrupting cells according to the invention has similar performance to the performance of the bench top bead beating method, which is well known in the art. In FIG. 7 , the horizontal axis represents the actuation time, i.e., the time taken to disrupt cells in the device, and the longitudinal axis represents Ct. An embodiment of the micro-device for disrupting cells according to the invention facilitates the disruption of cell membranes or cell walls by beating with microbeads. Therefore, the efficiency of elution of a particular substance in a cell, such as a nucleic acid, substantially increases, while the time for preparation of a sample, such as a nucleic acid preparation for a diagnostic assay, and the cost of the diagnostic assay are substantially reduced. In such an embodiment, the micro-device may be used in various diagnostic tools, such as a PCR apparatus, a microarray apparatus and a sequencing apparatus, for example, thereby increasing the accuracy of the diagnosis performed by the various diagnostic tools. Example 2: Cell lysis effect with respect to liquid volume fraction 1. Manufacture of Cell or Virus Lysis Device A cell or virus lysis device may be manufactured by disposing a commercially available elastic membrane between two glass chips. Chambers or channels may be defined using an entirety of a glass chip or a portion of the glass chips, which may then be combined with the elastic membrane therebetween, thereby completing the manufacture of the cell or virus lysis device. In the current example, first and second glass chips were manufactured by defining channels and chambers in a glass wafer by photolithography, etching, and wet etching processes, which are well known and the details of which can be determined without undue experimentation. After having been cleaned with a Piranha solution (i.e., a combination of sulfuric acid and hydrogen peroxide), a 6-inch glass wafer (borosilicate, 700 μm thick) was deposited with a 500 nanometer (nm) thick amorphous silicon layer by low-pressure chemical vapor deposition (“LPCVD”). Then, a patterning process was performed on a portion of the deposited silicon layer exposed through a photoresist film. The exposed part of the silicon layer was removed by dry etching. Afterwards, the photoresist film was stripped off, and the exposed glass wafer was wet-etched with a hydrofluoric acid solution (HF, 49%) to form a channel having a depth of about 100 μm and a width of about 200 μm. In the etching process for forming the channel, a weir (protrusion) projecting about 20 μm toward an inner center from an inner surface of the channel was formed to confine beads. The weir was formed to serve both as a valve seat and a bead trapping weir. Then, after the silicon layer was removed, a dry film resist (“DFR”) was coated and patterned. Then, a chamber including beads (ca. about 15.5 μl) and holes for fluid inflow or outflow were formed using a sand-blasting method. Subsequently, the glass wafer was diced into chip-shaped pieces, which were then washed with plasma. A fluidic chip (“first glass chip”) including the above-manufactured chamber to contain the beads, and a pneumatic chip (“second glass chip”) including a chamber to function as a pneumatic pump and not containing beads were permanently coupled with a 254 μm thick PDMS layer (available from Rogers Corporation) which was activated with a plasma between the first and second glass chips as an intermediate layer. The PDMS layer, which is a monolithic flexible layer, was used to control fluid flow and used as a pump and a valve and as an actuator for inducing collisions between the beads by pneumatic vibration. About 15-16 milligrams (mg) (about 2×10 5 in number) surface-modified glass beads were put into the bead chamber, which was then sealed using a PCR-compatible adhesive tape (available from Applied biosystems). The attached tape was covered with a polycarbonate plate to prevent the tape from bending during operation such as a DNA extraction. Operation of the PDMS layer was controlled by applying a positive pressure or a negative pressure to the pneumatic chamber with a Solenoid valve array (5070-5DC, available from SMC) connected thereto. The valves were coupled to an electropneumatic-regulator (ITV0030-3BL, available from SMC) and LabVIEW software (available from National Instruments). Operation of the valves associated with fluid transfer was visualized through an interface of the LabVIEW software in each step to monitor extraction of nucleic acids. FIGS. 8 and 9 illustrate a cell lysis device used in the current Example. FIG. 8 is a cross-sectional view of the cell lysis device using bead beating in which collisions of beads are induced by the vibration of a PDMS layer. Referring to FIG. 8 , an inlet with a first protrusion 40 and an outlet with a second protrusion 42 are connected to an inlet port and an outlet port, respectively, via fluid channels 48 and 50 . First and second valve sheets 44 and 46 are formed on an upper plate, which defines the inlet and the outlet to correspond to the inlet and outlet ports, respectively. Pneumatic chambers, specifically first to fourth chambers 24 A to 24 D, are disposed in a lower plate and the first and the fourth chambers 24 A and 24 D correspond to the first and second valve sheets 44 and 46 , respectively. First to fourth ports 34 A to 34 D are fluidly connected to the first to fourth chambers 24 A to 24 D. FIG. 9 is an enlarged view of a 3-layered cell lysis device including monolithic glass, PDMS and glass, and fluidic and pneumatic components of the cell lysis device. 2. Glass Bead Modification After having been washed with a Piranha solution and then with distilled water, glass beads having a diameter of about 30 μm to 50 μm (available from Polysciences, Inc.) were filtered and vacuum-dried. Afterwards, a 5% (volume/volume) trimethoxysilylpropyl-modified PEI (poly(ethyleneimine)-trimethoxysilylpropyl: PEIM) solution (available from Gelest, Inc.) in ethanol was prepared as a bead-surface modification solution. The beads were put into the bead-surface modification solution and reacted for about 2 hours by gentle mixing, followed by filtration and washing with fresh ethanol three times. The final recovered glass beads were incubated in a 110° C. oven for 50 minutes, to obtain glass beads having surfaces coated with PEIM. PEIM is known to be able to non-specifically bind to cells, and thus the glass beads coated with the PEIM may be used to nonspecifically separate cells. 3. Nucleic Acid Extraction A 1 mL sodium acetate buffer (50 millimolar (mM), pH 4, available from Sigma-Aldrich) containing 10 6 colony-forming units per milliliter (CFU/mL) of a sample S. aureus or Methicillin-resistant Staphylococcus aureus (“MRSA”), a 0.5-mL tris(hydroxymethyl)aminomethane (“Tris”)-ethylenediaminetetraacetic acid (“EDTA”) (“TE”) buffer (10 mM Tris, 1 mM EDTA, pH 8.0, available from Ambion) for washing, and a 10 μl or 20 μl NaOH (0.02 normal (N), available from Sigma-Aldrich) for lysis were stored into a liquid reservoir beforehand. The liquid solution was transferred by a pressure-driven operation. An operating liquid pressure was determined through a preliminary test. Initially, while applying a pressure of 150 kiloPascals (kPa) from above the PDMS layer, the sample solution was directed through the chamber containing the beads at 30 kPa and at about 200 μl/minute. After flowing through the chamber, the solution was recovered to evaluate a cell capture efficiency. After the initial loading of the sample, the washing solution was directed through the chamber containing the beads at about 500 μl/min (80 kPa) to wash it, which was then air-dried at about 100 kPa for about 30 seconds. To lyse the captured cells, 6 μl of the lysis solution (0.02 N NaOH) was injected into the chamber containing the beads, and valves on opposite sides of the chamber were closed. Subsequently, pressures of two pneumatic chambers were adjusted to about 80 kPa and about −80 kPa, respectively, by an adjustment of a Solenoid valve controlled via a Labview program, to vibrate the PDMS layer at a frequency of about 10 Hz for about 5 minutes, thereby performing a cell lysis process. After the cell lysis process, with the inlet and the outlet opened, 4 μl or 14 μl of the NaOH solution was injected with an application of a fluid pressure of about 100 kPa, to recover a lysed cell product through the outlet. The resulting lysed cell product containing a nucleic acid was 10 μl or 20 μl in total. The overall process took about 20 minutes or less. No additional DNA purification was performed. Experiments with a positive lysis control (“PLC”) and a negative lysis control (“NLC”) were conducted as follows. The experiment with a PLC was conducted using two different top bench lysis methods: an enzymatic method and a bead beating method. Two 1 mL sample solutions containing 10 4 CFU/mL and 10 6 CFU/mL of S. aureus were centrifuged in microcentrifuge tubes at about 13, 200 revolutions per minute (rpm) for about 20 minutes to precipitate cells. Then, the supernatant was removed from each centrifuged product. The precipitated pellets were treated using the two methods. In the enzymatic method, after an incubation of the cell pellets with a lysostaphin solution (200 milligrams per milliliter (mg/mL), available from Sigma) at about 37° C. for about 30 minutes, 20 μl of a purified DNA solution was obtained from the incubated product by using a Qiagen DNA extraction kit (Cat 51304, QIAamp DNA Mini Kit) according to the operation protocols of the kit. In the bead beating method, after an addition of 30 mg of bare glass beads and 20 mL of the lysis solution (0.02 N NaOH solution or distilled water) to the cell pellets, the combination was vigorously vortexed using a vortexer (GENIE 2, available from Fisher Scientific) at a full velocity for about 5 minutes. After simple centrifugation, an extracted DNA solution (a lysed cell product) was recovered. In the experiment with the NLC, after an addition of distilled water alone, the cell pellets were vigorously vortexed with no glass beads present, using a vortexer (GENIE 2, available from Fisher Scientific) at a full velocity of about 3,200 rpm for about 5 minutes. Using the resulting lysed control product, PLC, and NLC as templates, a target nucleotide sequence was amplified. This amplified product was compared with a result of amplifying the target nucleotide sequence with the DNA solution extracted using the bead beating based cell lysis device used as a template. For accurate comparison, a total number of S. aureus cells injected into the chamber and a final volume of the DNA extraction solution in each control group were controlled to be consistent with those of the test sample. 4. Real-Time PCR Amplification To lyse cells and qualify extracted DNA, real time-PCR was conducted using a GenSpector R TMC-1000 instrument (available from Samsung Electronics). A primer set (forward: 5′-GTT GCA TCG GAA ACA TTG TGT-3 (SEQ ID No. 1), reverse: 5′-ATG ACC AGC TTC GGT ACT ACT AAA GAT-3′ (SEQ ID No. 2), and GeneBank accession number AF033191) specific to the SA442 fragment of the S. aureus genome, and a primer set (forward: 5′-ACG AGT AGA TGC TCA ATA-3′ (SEQ ID No. 3), reverse: 5′-GGA ATA ATG ACG CTA TGA T-3′ (SEQ ID No. 4), and GeneBank accession number EF190335.1) specific to the mecA fragment of the MRSA genome were designed using a Primer3 software (Whitehead Institute/MT Center for Genome Research). A PCR reaction mixture (about 2 μl) was prepared to have the following concentration: 0.4 μM of Taqman probe (FAM-5′-TGT ATG TAA AAG CCG TCT TG-3′-MGB-NFQ (SEQ ID No. 5) for S. aureus ; and FAM-5′-CCA ATC TAA CTT CCA CAT ACC ATC T-3′-BHQ1 (SEQ ID No. 6) for MRSA), a 1× Z-Tag buffer (available from Takara Bio), 1 micromolar (μM) of each primer (available from Applied Biosystems or Sigma), 0.05 U of Z-Taq polymerase (available from Takara Bio), a 0.2 mM dNTP (available from Takara Bio), 0.5 μl of PCR-grade water (available from Ambion), and 1 μl of an extracted DNA solution. After the PCR reaction mixture was loaded into a PCR chip, thermal cycling was conducted as follows: denaturation at 95° C. for 1 minute and elongation at 60° C. for 4 seconds. The PCR conditions were designed to attain a PCR amplicon size of 72 base pairs (bp) for S. aureus and 98 bp for MRSA. The PCR amplicon sizes were further identified by gel electrophoresis (Agilent 2100 Bioanalyser, available from Agilent Technologies). 5. Experimental Results (1) Experimental Results of Control Sample Cell lysis and/or DNA purification effects were evaluated using a threshold cycle (Ct). Ct values of PLC samples are shown in Table 1. The Ct values in Table 1 are an average from three repeated experiments on each group. TABLE 1 Number of loaded S. aureus cells Benchtop bead beating method Enzymatic method (CFU) NaOH (0.02 N) Distilled water (lysostaphin) About 10 4 30.5 ± 0.35 34.5 ± 0.26 31.6 ± 0.56 About 10 6 23.7 ± 0.29 26.7 ± 0.15 25.7 ± 1.43 To improve lysis efficiency, the benchtop bead beating method was performed with a lysis solution containing a surfactant or a chemical substance. In the current example, a NaOH solution (0.02 N), known not to interfere with PCR amplification without additional purification, was used. Bead beating effects with the NaOH solution or distilled water in Table 1 indicate that NaOH, which chemically destructs cell walls, contributes to improving the DNA extraction efficiency. While not wanting to be bound by theory, regarding the enzymatic lysis of S. aureus , lysostaphin was selected because it specifically cleavages cross-linking pentaglycine bridges in the cell wall of staphylococci. The benchtop bead beating method with the NaOH solution shows a performance more than or equivalent to that of the enzyme-based DNA extraction method. Therefore, results with the used cell or virus lysis device were compared with those with the benchtop vortexing machine to evaluate efficiency. Optical density measurement is an approximate quantification of cells, which have caused variations in the Ct value of the PLC even at the same optical density with a standard deviation of about 1.5. The NLC sample vortexed with distilled water alone had a Ct value of about 31.5 at a cell concentration of 10 6 CFU/mL. (2) Cell Capture Results In the present example, basic operations were as follows: (1) capturing cells on glass beads, (2) washing and drying, (3) lysing cells with in-situ bead beating, and (4) eluting the extracted DNA solution. Bacteria cells may be specifically or non-specifically captured on a solid substrate. In addition to cell specific immunoaffinity techniques, pathogenic bacteria cells may be captured by non-specific cell capture techniques using surface thermodynamics or electrostatic interaction. In the present example, to evaluate the effect of long-range Coulombic electrostatic interactions to non-specific cell capture glass beads were modified to be positively charged on the surface. Surfaces of the glass beads were treated to have positive amine derivatives by reaction with an organosilane compound including poly(ethyleneimine) (“PEI”). After the interior of the microfluidic chamber was filled with the surfacemodified glass beads, 1 mL of a sample solution containing 10 6 CFU/mL of S. aureus cells was directed through the microfluidic chamber. Initially, results with non-modified glass beads and those with modified glass beads were compared in terms of electrostatic interaction. As a result, a Ct value of the DNA extracted with the modified glass beads was smaller by 2 than that of the DNA extracted with the non-modified glass beads, thereby indicating an increase of about 4 times in cell capture efficiency through the surface modification, considering that a difference in Ct of one (1 Ct) indicates two times the difference in initial template copy number. To obtain quantitative data of the cell capture capacity of the cell lysis device used in the present example, after flowing through the chamber, the sample solution was recovered, centrifuged, and disrupted using a benchtop vortexing machine together with bare glass beads and an NaOH solution, as used with the PLC (hereinafter, the lysed product is referred to as “AC”). A Ct value of the AC was compared with that of a DNA solution (“AE”) eluted from the bead-beating device through appropriate operations. A difference between the two values was used as a measure of cell capture efficiency and capacity of the bead beating device. When 1 mL of each S. aureus sample solution having a different cell concentration of 10 4 , 10 5 , 10 6 , 10 7 , and 10 8 CFU/mL was loaded, the differences between the AC and AE were maintained at about 5 or greater in the sample solutions of from 10 3 CFU/mL to 10 7 CFU/mL, but was reduced to about 2.5 at 10 8 CFU/mL. Considering that one hundred times the difference in initial template copy number may induce a difference in Ct (ΔCt) of about 3.3, the capture efficiency was found to be about 90%. The manufactured bead-packed microfluidic device had a capacity of 10 7 CFU or greater of S. aureus cells. After cell capture, the bead-packed microchamber was washed, dried with air, and then filled with the lysis solution. These results indicate that the applied cell capture method is appropriate to disrupt the captured cells by in-situ bead beating, rather than to release the captured cells. Other cell capture methods using, for example immunoaffinity, may be integrated with the cell lysis device used in the present example by employing appropriate solid surface chemistry. (3) Effects of Liquid Volume Fraction (f L ) on Bead Beating Cell Lysis Effects of various factors on bead beating cell lysis were investigated. As a result of experiments at a membrane vibration frequency (from about 5 Hz to about 10 Hz), a membrane operating pressure (from about 20 kPa to about 80 kPa), and a depth of a pneumatic displacement chamber (from about 100 μm to about 200 μm), the factors were found not to be statistically significant on PCR Ct values. Effects of liquid viscosity on bead beating cell lysis were investigated with respect to changes in liquid volume fraction (f L ) during bead beating induction, the liquid volume fraction being defined by the volume of the lysis solution with respect to the pure void volume of the microfluidic chamber packed with beads. The final elution volume of DNA was adjusted to be 20 μl, by adding a NaOH solution when the extracted DNA was eluted from the chamber. FIG. 10 is a graph of threshold cycle (Ct) and change in Ct (ΔCt) versus liquid volume fraction (f L ) showing cell lysis efficiency with respect to changes in liquid volume fraction (f L ). Referring to FIG. 10 , the liquid volume fraction (f L ) is found to correlate with cell lysis efficiency, and, while not wanting to be bound by theory, it is believed that liquid volume fraction (f L ) is a determining factor of cell lysis efficiency. The cell lysis efficiency was high at a liquid volume fraction (f L ) of 0.6 or less, and in particular, was higher at 0.5 or less, and still higher at between 0.3 and 0.5. Efficient cell lysis was also possible at a liquid volume fraction (f L ) of 0. This is attributed to the fact that the viscosity of the liquid solution (NaOH) is 100 times or greater than that of air, and thus, a relative amount of the liquid solution (NaOH) markedly affects the viscosity of the mixed solution (gas and liquid). It should be understood that the embodiments described herein should be considered in a descriptive sense only and not for purposes of limitation. Descriptions of features or aspects within each embodiment should typically be considered as available for other similar features or aspects in other embodiments.

A micro-device for disrupting cells includes a first chamber in which the cells are disrupted, a second chamber which is pressurized and depressurized, a flexible membrane which separates the first chamber and the second chamber and is vibrated by pressuring and depressurizing the second chamber, and a micro-unit confined in the first chamber, where the micro-unit disrupts the cells in the first chamber.

This application claims priority under 35 U.S.C. §§119 and/or 365 of Appln. No. 100 40 869.9 filed in Germany on Aug. 21, 2000, the entire content of which is hereby incorporated by reference. FIELD OF THE INVENTION The invention relates to a method for the fluid-mechanical stabilization of a premix burner, into which a combustion air stream is fed tangentially into an interior burner chamber, is mixed with an injected gaseous and/or liquid fuel while forming a coaxially oriented swirl flow and induces a reverse flow zone at a burner outlet that is used during the operation of the burner to stabilize the flame. The invention furthermore relates to a premix burner for performing the method. A preferred field of application of the invention is the operation of a gas turbine system. BACKGROUND OF THE INVENTION Premix burners of the type discussed here are known from EP 0 321 809 and EP 0 780 629. Such burners, characterized by very low noxious emissions, are used widely in combustors of gas turbine systems for hot gas generation. When operating gas turbine systems, thermoacoustic oscillations often occur in the combustors. Fluid-mechanical instability waves created at the burner result in the formation of flow vortices that have a major effect on the entire combustion process and result in undesirable, periodic heat releases within the combustor that are associated with major fluctuations in pressure. The high fluctuations in pressure are coupled with high oscillation amplitudes that can lead to undesirable effects, such as, for example, a high mechanical load on the combustor housing, increased NO x emissions caused by inhomogeneous combustion, and even an extinction of the flame within the combustor. Thermoacoustic oscillations are based at least in part on flow instabilities of the burner flow that express themselves as coherent flow structures and influence the mixing processes between air and fuel. In standard combustors, cooling air is passed in the form of a cooling air film over the combustor walls. In addition to the cooling effect, the cooling air film also has a sound-dampening effect and helps to reduce thermoacoustic oscillations. In modern high-efficiency gas turbine combustors with low emissions and constant temperature distribution at the turbine inlet, the cooling air flow into the combustor is clearly reduced, and the entire air is passed through the burner. However, at the same time the sound-dampening cooling air film is reduced, causing a reduction in the sound-dampening effect so that there is once again an increase in the problems associated with undesirable oscillations. Another possibility for dampening the sound is the connection of so-called Helmholtz resonators near the combustor or cooling air supply. However, because of tight space conditions, it is very difficult to provide such Helmholtz resonators in modern combustion chamber designs. It is also known that the fluidic instabilities and associated pressure fluctuations occurring in the burner can be countered by stabilizing the fuel flame with an additional injection of fuel. Such an injection of additional fuel is performed through the head stage of the burner that is provided with a jet for the pilot fuel gas supply located on the burner axis; however, this results in an over-rich central flame stabilization zone. This method of reducing thermoacoustic oscillation amplitudes has the disadvantage, however, that the injection of fuel at the head stage may be associated with increased NO x emissions. Closer studies regarding the formation of thermoacoustic oscillations have shown that such undesirable, coherent structures are formed during mixing processes. Of special significance hereby are the shear layers forming between two mixing flows, which are formed within the coherent structures. More detailed explanations regarding this phenomenon can be found in the following publications: Oster & Wygnansky, 1982, “The forced mixing layer between parallel streams,” Journal of Fluid Mechanics, Vol. 123, 91-130; Paschereit et al., 1995, “Experimental investigation of subharmonic resonance in an axisymmetric jet,” Journal of Fluid Mechanics, Vol. 283, 365-407. As can be seen from the previous articles, it is possible to influence the coherent structures forming within the shear layers by introducing an acoustic excitation in a targeted manner in such a way that their formation is prevented. Another method is the introduction of an acoustic anti-sound field so that the existing, undesired sound field is properly extinguished by a phase-shifted sound field that has been introduced in a targeted manner. The anti-sound technique, as described, requires a relatively high amount of energy, however, which must be provided to the burner system either externally or must be branched off from the entire system at a different point, which, however, will result in a small, yet existing, loss of efficiency. In addition to the previous active methods for specifically reducing the coherent structures forming inside burners, such interferences in the burner flow alternatively can be countered with passive measures. Passive measures, i.e., primarily constructive design characteristics of the burners, that extend the operating range of a burner with respect to pulsations and emissions are especially attractive, since, once installed, they do not require any additional maintenance. SUMMARY OF THE INVENTION The invention is based on the objective of creating a method for increasing the fluidic stability of a premix burner, which efficiently and without further energy consumption suppresses the undesired flow eddies that form as coherent pressure fluctuation structures. The measures necessary on a premix burner for this purpose should be simple to construct and cheap to realize. The measures used also should be completely maintenance-free. According to the invention, the objective is realized with a method for increasing the fluidic stability of a premix burner as well as with a premix burner of the type mentioned in the independent claims. Characteristics that constitute advantageous further development of the concept of the invention are described in the dependent claims and the specification, as well as in the exemplary embodiments. The method according to the invention is based on the basic idea of—for the fluid-mechanical stabilization of a premix burner, into which at least one combustion air stream is fed tangentially into a burner chamber and is mixed with an injected gaseous and/or liquid fuel while forming a swirl flow oriented coaxially to the burner axis and induces a reverse flow zone at a change in the cross-section on a burner mouth, that is used during the operation of the burner to stabilize the flame—increasingly, radially deforming the swirl flow within the burner chamber in the direction of the burner mouth towards at least one circumferential section and let it enter the combustor in a non-rotation-symmetrical flow cross-section, whereby this deformation is created by reducing the free flow cross-section of the burner chamber. The formation of coherent vortex structures is hindered by a shape of the flow cross-section that deviates from the rotation symmetry in the burner chamber and on entering the combustor. In premix burners according to the state of the art, the time delay of the fuel from the injection point to the flame is constant at certain operating points. The deformation of the flow-cross-section according to the invention results in a broad distribution of the delay time. The prevention of the formation of vortex structures at the burner outlet and a smudged time delay also suppresses a periodic heat release, which again is responsible for the occurrence of thermoacoustic oscillations. By forcing the deformation of the swirl flow through constricting sections of the chamber contour, as will still be explained at another place, this results in an acceleration of the flow, which acts in a stabilizing manner on the reverse flow zone. A premix burner according to the invention is based on a premix burner for use in a heat generator, comprising essentially a swirl flow generator with means for the tangential introduction of at least one gaseous and/or liquid fuel into the combustion air stream with concomitant formation of a swirl flow with an axial movement component up to the burner mouth, at which the swirl flow bursts while inducing a reverse flow zone. A burner according to this type, based on at least two hollow, conically expanding partial bodies stacked inside each other in the flow direction of the hot gases, the center axes of which extend offset to each other, is described in EP 0321809, which is an integrated part of this application. Such burner types, also called cone burners or double cone burners, are provided at their burner outlet with a separation edge, whose edge contour consists of two semi-circles, offset from each other, but whose edge contour is almost circular and therefore approximately rotation-symmetrical to the burner axis when closed. The fuel/gas mixture forming inside the burner chamber spreads in the form of a rotation-symmetrical swirl flow with an axial component towards the burner mouth, with all its known disadvantages with respect to the formation of coherent structures and associated thermoacoustic pressure fluctuations. If, however, it is ensured that targeted radial deformations are introduced into the flow of the fuel/air mixture in such a way that the flow cross-section differs from that of an rotation-symmetrical flow, the formation of coherent structures and a constant time-delay of the fuel can be effectively countered in this way. Such an influencing of the flow geometry can be achieved with at least one section of the chamber wall, where said wall section has a smaller slope in a down-stream end part of the burner chamber than in an upstream part. In this way, this at least one section, in contrast to those wall sections at the same axis level that do not possess this property, results in a radial deviation from the circular shape in the direction towards the burner axis. Any partial non-circular contours in the burner chamber and at the outlet edge, for example, straight or spherically curved edge sections along the circumference, help in reducing flow vortices. It should be observed as a principal design rule for designing the burner outlet edge that the geometrical deviation from a circular geometry should be chosen at least so large that the resulting distance between the two geometries is greater than the boundary layer thickness of the flow that flows through the outlet geometry. BRIEF DESCRIPTION OF THE DRAWINGS The invention is described in an exemplary manner below with the help of exemplary embodiments in reference to the drawing, without limiting the general concept of the invention whereby: FIG. 1 a shows a perspective drawing of a premix burner according to the state of the art, on which the invention is based; FIG. 1 b shows another drawing of a burner in a simplified form; FIGS. 2 a-d show a very schematic portrayal of the concept of the invention using various forms of swirl flow generators; FIG. 3 shows an embodiment of burner modified according to the invention; FIG. 4 shows a burner according to an embodiment of the invention having sections of the chamber wall that constrict the flow cross-section; FIGS. 5 a and 5 b show an axial cross-section and an end view of a premix burner according to an embodiment of the invention having a mixing section that is rotationally symmetrical about a central axis; FIGS. 6 a and 6 b show an axial cross-section and an end view of a premix burner according to an embodiment of the invention having a mixing section that is rotationally asymmetrical about a central axis; FIGS. 7 a and 7 b show an axial cross-section and an end view of a premix burner according to an embodiment of the invention having a cylindrical or convergent nozzle section at the downstream burner end; FIG. 8 shows a portrayal of the suppression of combustion oscillations by suppressing flow vortices in a burner. DETAILED DESCRIPTION OF THE INVENTION FIGS. 1 a and 1 b show in a very schematic form the construction and function of a premix burner that is the starting point of the present invention. The premix burner comprises two hollow, conically expanding partial bodies ( 1 ) and ( 2 ) arranged axis-parallel and offset relative to each other in such a way that they form tangential slits ( 3 ) in two overlapping areas located in a mirror-image opposite from each other. Although FIGS. 1 a and 1 b show, as an example, two conically expanding partial bodies ( 1 ) and ( 2 ), other configurations are also conceivable. These burners, for example, are not limited to the arrangement of two partial bodies ( 1 ) and ( 2 ), nor is their conical configuration obligatory. The expert will be aware of this. The gaps ( 3 ) resulting from the offset of the longitudinal axes are used as inlet channels through which the combustion air ( 5 ) flows tangentially into the burner chamber ( 6 ) during the operation of the burner. Injection openings ( 7 ) through which a preferably gaseous fuel is injected into the combustion air ( 5 ) that is flowing by are located along the tangential inlet channels ( 3 ). In the interest of good mixing, the fuel injection takes place preferably within the tangential inlet channel ( 3 ), immediately before the entrance into the burner chamber ( 6 ). The beginning part of the burner, which also may be constructed cylindrically (not shown), a central nozzle ( 8 ) for atomizing a liquid fuel is provided, the capacity and function of which nozzle depends on the burner parameters. The liquid fuel leaves the nozzle ( 8 ) at an acute angle and forms a cone-shaped fuel profile in the burner chamber ( 6 ), which fuel profile is enclosed and continuously broken down into a mixture by the tangentially entering combustion air ( 5 ) that changes into a swirl flow ( 9 ), which process can be supported by preheated combustion air ( 5 ) or by mixing in recycled waste gas. Alternatively, it is also possible to supply the nozzle ( 8 ) with gaseous fuel. On the Combustor side, the premix burner has a front palate ( 10 ) functioning as an anchor for the partial bodies ( 1 ) and ( 2 ), which is provided with a number of drilled holes ( 11 ) for introducing air into the combustor ( 12 ). The fuel/air mixture passing through the burner chamber ( 6 ) in a swirl flow ( 9 ) reaches the optimum fuel concentration across the cross-section at the downstream end of the premixing section ( 13 ) at the burner mouth ( 14 ). When exiting from the burner, the swirl flow ( 9 ) bursts, forming a reverse flow zone ( 15 ) with a stabilizing effect for the flame front ( 17 ) functioning there. This aerodynamic flame stabilization quasi assumed the function of a flame holder. The feared failure of mechanical flame holders due to overheating, followed possibly by serious failures of machine sets, is therefore prevented. In addition, the flame does not lose any heat to the cold walls, except by radiation. This also aids in homogenizing the flame temperature and therefore contributes to lower noxious emissions and good combustion stability. According to the invention, measures are now provided to increasingly deform the swirl flow ( 9 ) inside the premix section ( 13 ) radially. It is preferred that this deformation takes place symmetrically. However, this is not mandatory. An important characteristic hereby is that this deformation is brought about by reducing the free flow cross-section ( 18 ). The wall ( 21 ) of the chamber ( 6 ) has in a downstream part ( 20 ) at least one section ( 22 ) that has a smaller slope with respect to the burner axis ( 4 ) than an upstream part ( 19 ). This means that the contour ( 21 ) of the burner chamber ( 6 )—which, when seen over its cross-section, is approximately circular—is provided with sections ( 22 ) that are distributed over the circumference and deviate from the circular shape of the chamber contour ( 21 ) towards the center axis ( 4 ), i.e., constrict the chamber ( 6 ), as is shown in the longitudinal section schematically shown in FIGS. 2 a - 2 d . In this context it was found to be advantageous, after all, that the deformation of the flow is simultaneously accompanied by an acceleration of the flow. This measure has a particularly beneficial effect on the stability of the burner. On the one hand, the cross-section shape of the flow ( 9 ) exiting from the burner, which shape deviates from the rotation symmetry, impairs the formation of coherent vortex structures and ultimately inhibits the generation of thermoacoustic oscillations thereby. On the other hand, the acceleration of the swirl flow ( 9 ) at the burner outlet ( 14 ) resulting from the absolute or relative constriction of the flow cross-section ( 18 ) brings about a stabilization of the reverse flow zone ( 15 ), inhibiting fluctuations of the reverse flow zone ( 15 ), the associated period heat release, and thus the development of thermoacoustic oscillations. This combination of equally acting effects results in synergy effects that permit, in a particularly advantageous manner, to increase the fluid-mechanical stability of a premix burner. FIGS. 2 a - 2 d are intended to explain the concept of the invention using very schematic drawings. FIG. 2 a shows a known swirl flow generator geometry that can be used to realize the invention in a particularly advantageous manner, whereby—as mentioned at another place—the conical configuration of the swirl flow generator ( 13 ) is not mandatory. FIGS. 2 b - 2 d symbolize the concept of the invention, which consists of angling the wall ( 21 ) of the burner chamber ( 6 ) in at least one circumferential section ( 22 ) by reducing the free flow cross-section ( 18 ) in the direction of the burner axis ( 4 ) in order to deform the flow profile. This may be accomplished symmetrically or asymmetrically with at least one such section ( 22 ) that constricts the flow cross-section. In a downstream part ( 20 ) of the chamber ( 6 ), which part ( 20 ) may start at approximately ⅔ of the axial length, the chamber wall ( 21 ) is bent in at least one circumferential section ( 22 ) at an angle in the range from 2° to 45°, in particular 5° to 15°, towards the burner axis ( 4 ). The expert also will be able to deduce from these schematic drawings another advantage of the invention, i.e., the possibility to retrofit existing burners with little expenditure. The sections ( 22 ) constricting the flow cross-section ( 18 ) can be realized with the help of flow-guiding installations ( 28 ) applied at a later time. FIGS. 3 to 7 show embodiments of burners designed according to the invention. FIG. 3 shows a preferred variation of the invention, according to which the burner mouth ( 14 ) has a polygonal outlet contour ( 16 ). As can be seen most clearly from the schematic drawings of FIGS. 2 a - 2 d , the conically expanding contour ( 23 ) of the burner chamber ( 6 ) is discontinued in a downstream end part ( 20 ) and is continued with a slope smaller than the previous part ( 19 ) in relation to the longitudinal axis ( 4 ). The term “smaller slope” also is supposed to include a progression parallel to the longitudinal axis ( 4 ) or a convergent progression, as shown in FIGS. 2 a - 2 d . The expert has a number of possibilities available to realize this suggestion. According to one preferred embodiment, appropriately shaped plates ( 28 ) are welded into the shell-shaped partial bodies ( 1 ) and ( 2 ) of burners constructed according to the state of the art, whereby these plates represent—seen two-dimensionally—chords that cut sectors from the free flow cross-section ( 18 ) of the burner chamber ( 6 ). It is preferred that for each partial body ( 1 ) or ( 2 ) preferably one to four such plates ( 28 ) are welded onto the inside wall ( 21 ). In new burners, the wall contour is shaped during the manufacturing process. According to another embodiment seen from FIG. 4 in connection with FIG. 2 c , the burner is constructed in an upstream part ( 19 ) in an actually known manner of two partial bodies ( 1 ) and ( 2 ) with an essentially circular cross-section that are stacked inside each other in an offset manner. In a transitional area, at approximately ⅔ of the axial length, the inside wall ( 21 ) changes from its essentially circular contour to a polygonal one that becomes increasingly more distinct in its further progression towards the burner mouth ( 14 ). These sections ( 22 ) of the chamber wall ( 21 ) that constrict the flow cross-section ( 18 ) like chords have less of a divergence in relation to the longitudinal axis ( 4 ) compared to the upstream parts ( 19 ) of the chamber wall ( 6 ). The term “less divergence” hereby shall also include the possibility of a parallel or convergent progression relative to the longitudinal axis ( 4 ). When viewing the cross-section, the constricting sections ( 22 ) as a rule have a linear contour. However, a slightly convex or concave progression is also possible. A convex progression is advantageous especially is only a small number or only one or two of such sections ( 22 ) are provided. Another embodiment, not shown in a figure, consists of not providing the burner chamber ( 6 ), even in its upstream part ( 19 ), with a circular cross-section, but to equip the burner with a chamber ( 6 ) with a continuously non-rotation-symmetrically contoured chamber ( 6 ). This embodiment is particularly suitable for polygonal contours ( 23 ) of the chamber cross-section ( 18 ). From the state of the art, it is known per se, to fit burners, as they were defined previously, for the purpose of better mixing and flame positioning for difficult fuels with nozzles ( 24 ) or mixing pipes ( 25 ) that follow the swirl flow generator ( 13 ). Even for these types of burner variations, the invention can be used to increase the fluid-mechanical stability of such burners by interfering with the flow instabilities and producing a smudged time delay of the fuel from the injection site to the flame. FIGS. 5 a , 5 b , 6 a and 6 b show a premix burner consisting of a swirl flow generator ( 13 ) for a combustion air stream ( 5 ) and means for injecting at least one fuel, whereby downstream from the swirl flow generator ( 13 ) a mixing section ( 25 ) is located. In the housing ( 26 ) surrounding the mixing section ( 25 ), inlet openings ( 27 ) for injecting an additional combustion air amount can be located evenly distributed over the circumference so as to extend at an acute angle to the longitudinal axis ( 4 ). It is preferred that in an area downstream from the inlet opening ( 27 ), the rotation-symmetrical flow cross-section of the mixing section ( 25 ) is deflected by sections ( 22 ) that construct the free circumference ( 29 ) and is radially deformed. The outlet opening ( 16 ) takes on a polygonal cross-section shape, composed of a plurality of linear sections ( 22 ). Very promising are outlet contours ( 16 ) in the form of a regular or irregular polygon. The individual, linear sections ( 22 ) of the outlet edge ( 27 ) span the outlet opening ( 16 ) of the burner. However, this linearity, as already mentioned before, is not mandatory, and these sections ( 22 ) also can be convex or concave. FIGS. 6 a and 6 b indicate a convex wall section ( 22 ) with an asymmetrical arrangement. FIGS. 7 a and 7 b , show an embodiment with a cylindrical or convergent nozzle section ( 24 ) at the downstream burner end. According to the state of the art, these downstream nozzles ( 24 ) primarily have the function of accelerating the flow at the burner outlet and thus stabilize the reverse flow zone ( 15 ). According to one embodiment of the invention, this desirable acceleration through a reduction in the cross-section that starts in flow direction and increases is achieved by constricting this nozzle section ( 24 ) in flow direction from an essentially circular cross-section shape to another cross-section shape, for example, a regular or irregular polygon or oval. FIG. 8 shows a diagram that shows the combustion power of the burner according to FIG. 3 along the abscissa, and a scale quantifying the formation of thermoacoustic oscillations created as a result of coherent structures within the flow inside the burner along the ordinate. The thermoacoustic oscillations shown are in the 100 Hz range. If a burner with conventional burner outlet according to the embodiment in FIG. 1 (line with squares) is compared with a burner outlet according to the invention as shown in the embodiment in FIG. 3 (line with circles), it is clear that in the latter significantly less thermoacoustic oscillations are created. The previously described embodiments should in no way be seen in a sense that would limit the invention. They are instructive and should be understood as an outline of the many possible embodiments of the invention as characterized in the claims. LIST OF REFERENCE SYMBOLS 1 partial body 2 partial body 3 tangential combustion air inlet channel 4 burner axis 5 combustion air 6 burner chamber 7 injection openings for fuel 8 central fuel nozzle 9 swirl flow 10 front plate 11 cooling air holes 12 combustor 13 swirl flow generator, premixing section 14 burner mouth 15 reverse flow zone 16 outlet cross-section into combustor 17 flame front 18 flow cross-section of burner chamber 19 upstream part of burner chamber 20 downstream part of burner chamber 21 wall of burner chamber 22 wall section 23 inside contour of burner chamber 24 burner nozzle 25 mixing section 26 mixing section housing 27 outlet edge 28 installations 29 flow cross-section in the mixing section 30 wall of mixing section 31 upstream part of mixing section 32 downstream part of mixing section 33 flow cross-section of burner nozzle 34 wall of nozzle 35 upstream part of nozzle 36 downstream part of nozzle

The subject of the invention is a method as well as an apparatus for suppressing flow vortices within a turbo power machine with a premix burner, into which fuel and air are introduced for mixing, which then leave the burner downstream along its burner axis in the form of a fuel/air mixture through a burner outlet and flow into a combustor located downstream from the burner in the flow direction of the fuel/air mixture. The invention is based on the basic idea of—for the fluid-mechanical stabilization of a premix burner, into which at least one combustion air stream ( 5 ) is fed tangentially into a burner chamber ( 6 ) and is mixed with an injected gaseous and/or liquid fuel ( 7;8 ) while forming a swirl flow ( 9 ) oriented coaxially to the burner axis and induces a reverse flow zone ( 15 ) at a change in the cross-section on a burner mouth ( 14 ), that is used during the operation of the burner to stabilize the flame—increasingly, radially deforming the swirl flow ( 9 ) within the burner chamber ( 6 ) in the direction of the burner mouth ( 14 ) and let it enter the combustor ( 12 ) in a non-rotation-symmetrical flow cross-section, whereby this deformation is created by reducing the free flow cross-section ( 18 ) of the burner chamber ( 6 ). The fuel/air mixture flows into the combustor with a non-rotation-symmetrical flow cross-section.

BACKGROUND OF THE INVENTION This invention relates to an improvement in vehicular automatic speed change gear assemblies. A conventional automatic speed change gear assembly is designed so that an oil pressure control circuit thereof includes a manual valve having the six positions P, R, N, D, 2 and L to permit the speed change gear mechanism to provide a plurality of speed ranges. Accordingly, the coventional automatic speed change gear assembly is disadvantageous in the following points: It is necessary to provide a large space in the vehicle in order to provide the above-described operating lever. Furthermore, the operator may operate the operating lever erroneously. More specifically, the manual valve may be set to the position P for parking the vehicle or to the position R for running the vehicle in reverse while the vehicle is moving forwardly, or the manual valve may be switched to the position N in returning it to position D from either second or low. Furthermore, in order to park the vehicle, it is necessary to shift the operating lever to a position corresponding to the position P of the manual valve. Accordingly, the conventional automatic speed change gear assembly suffers from difficulties in that it is troublesome to shift the operating lever to the P position whenever it is required to park the vehicle, and if the operator forgets to so set the operating lever in parking the vehicle on a slope, a hazard occurs in that the vehicle may roll unintentionally. SUMMARY OF THE INVENTION In view of the foregoing, a primary object of this invention is to provide an automatic speed change gear assembly in which the above-described hazards accompanying the erroneous operation of a conventional automatic speed change gear assembly are eliminated, and wherein operation is carried out completely automatically, and is simplified. Another object of the invention is to provide an automatic speed change gear assembly in which the operating conditions of the vehicle are detected, so that the parking mechanism of the vehicle is automatically engaged, whereby the operability and the security thereof are improved. The foregoing and other objects of the invention have been achieved by the provision of an automatic speed change gear assembly designed so that only when a signal from a vehicle stop sensor for detecting the stop of the vehicle is available are forward and backward movement instructions from the operator made effective. Furthermore, the objects of the invention have been achieved by the provision of an automatic speed change gear assembly designed so that when the signal from the vehicle stop sensor and an engine "off" signal from a key switch controlling the start and stop of the engine are available, the parking mechanism actuator of the vehicle is placed in an engagement state. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is a power train chart showing an automatic speed change gear assembly according to one embodiment of this invention; FIG. 2 is a table indicating the relationships between the operations of the frictional elements of the gear train and the speed ranges of the automatic speed change gear assembly; FIGS. 3A and 3B are a circuit diagram showing an oil pressure control device of the automatic speed change gear assembly; FIG. 4 is a table indicating the relationships between solenoid valve operations and the speed ranges; and FIGS. 5A, 5B, and 5C are a flow chart describing the control of an operating condition change-over valve and the solenoid valves of the oil pressure control device. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS One embodiment of this invention will now be described in detail with reference to the accompanying drawings. As shown in FIG. 1, the crank shaft 4 of an engine 2, which is the power source for the vehicle, is connected directly to the pump 8 of a torque converter 6. The torque converter 6 comprises the pump 8, a turbine 10, a stator 12 and a one-way clutch 14. The stator 12 is coupled through the one-way clutch 14 to a casing 16. By means of the one-way clutch 14, the stator 12 is allowed to turn in the same direction as the crank shaft 4, but it is not allowed to turn in the opposite direction. The torque transmitted to the turbine 10 is transmitted through an input shaft 20 to a speed change gear assembly 22 which provides four forward ranges and one reverse range. The speed change gear assembly 22 comprises three clutches 24, 26 and 28, two brakes 30 and 32, a one-way clutch 34, and a planet gear mechanism 36. The planet gear mechanism 36 is made up of a ring gear 38, a long pinion gear 40, a short pinion gear 42, a front sun gear 44, a rear sun gear 46 and a rotatable carrier 48 which rotatably supports the pinion gears 40 and 42. The ring gear 38 is coupled to an output shaft 50. The front sun gear 44 is coupled through a kick down drum 52 and the front clutch 24 to the input gear 20. The rear sun gear 46 is coupled through the rear clutch 26 to the input shaft 20. The carrier 48 is coupled to the casing 16 through the low reverse brake 32 and the one-way clutch 34 which are arranged functionally in parallel. The carrier 48 is further coupled to the input shaft 20 through the fourspeed clutch 28 provided at the rear end of the speed change gear assembly 22. The aforementioned kick down drum 52 can be fixedly coupled to the casing 16 through the kick down brake 30. Engaging teeth 54 formed on the periphery of the ring gear 38 and a pawl 58 form a parking mechanism 59. The pawl 58 is extended by means of an actuator 56 including a motor and a solenoid, which is operated by a signal from an electronic control device (described later), as a result of which the pawl 58 is engaged with the engaging teeth 54, to thereby prevent rotation of the ring gear 38 (and accordingly the output shaft 50). The torque passing through the planet gear mechanism 36 is transmitted through an output gear 60 secured to the output shaft 50 and an idler gear 62 to a driven gear 64. Torque is further transmitted through a transfer shaft 66 secured to the driven gear 64 and a helical gear 68 to a differential gear device 72 to which a drive shaft 70 is coupled. The above-described clutches and brakes are frictional engagement devices including engaging piston units or servo units, and are operated by the oil pressure produced by an oil pump 74 driven by the engine 2 when coupled to the pump 8 of the torque converter 6. The oil pressure is selectively applied to the clutches and the brakes by an oil pressure control device (described hereafter) according to operating conditions detected by various operating condition detecting devices. The combinations of the operations of the clutches and the brakes provide the speed change ranges, namely, four forward ranges and one reverse range. FIG. 2 shows the relationships between the operations of the clutches and brakes and the various speed ranges. In FIG. 2, "O" indicates that the clutch or brake is in engagement, and " " indicates that, immediately before the low reverse brake 32 engages to effect a speed change, the rotation of the carrier 48 is stopped by the action of the one-way clutch 34. The oil pressure and electronic control systems allowing the speed change gear assembly 22 (in FIG. 1) to provide the speed change ranges shown in FIG. 2 will now be described. The oil pressure control device as shown in FIGS. 3A and 3B operates to selectively supply the pressurized oil, which is delivered from an oil pool 76 through an oil filter 78 and an oil path 80, according to the operating conditions to operate the piston units or servo units of the clutches 24, 26 and 28 and the brakes 30 and 32 of the speed change gear assembly 22 and the torque converter 6. The oil pressure control device comprises a vent valve 82, a torque converter control valve 84, a reducing valve 86, an operating condition change-over valve 88, a shift control valve 90, a rear clutch control valve 92, an N-R control valve 94, an oil pressure control valve 96 operated during a speed change, an N-D control valve 98, a first-second speed shift valve 100, a second-third speed and fourth-third speed shift valve 102, a four-speed clutch control valve 104, and solenoid valves 106, 108 and 110. These circuit elements are variously connected through oil paths to one another. The solenoid valves 106, 108 and 110 are all of the same construction. More specifically, the solenoid valves are of the normally closed type which open and close orifices 114, 116 and 118 in response to electrical signals from an electronic control device 112, respectively. The solenoid valves 106, 108 and 110 have solenoids 120, 122 and 124, valves 126, 128 and 130 arranged in the solenoids to open and close the orifices 114, 116 and 118, and springs 132, 134 and 136 urging the valves 126, 128 and 130 to their closed positions, respectively. The electronic control device 112 incorporates an operating condition determining device for detecting the operating conditions of the vehicle to determine the position of the operating condition change-over valve 88 and the open and close combinations of the solenoid of the solenoid valves 108 and 110 according to a process described hereafter, and a speed change detecting device for detecting the beginning of a speed change. The electronic control device 112 operates to control the position of the operating condition change-over valve 88, actuate and stop the solenoid valve 106 (which is subjected to duty control), control the pulse width of a pulse current at 50 Hz supplied to the solenoid valve 106 to change the valve opening time thereof to thereby control the oil pressure, and open and close the solenoid valves 108 and 110. Inputted to the electronic control device are signals from an advancement direction indicating switch 138 (hereinafter referred to merely as "an indicating switch 138") which is provided in the vehicle (not shown) where it can be operated readily by the operator and which indicates backward movement when turned on and forward movement when turned off; a door switch 140 for detecting the operation of the door by the driver's seat, the door switch 140 being turned on when the door is closed; a seat switch 142 provided on the driver's seat, the switch 142 being turned on when the driver sits on the seat; a side switch 144 for detecting the operating condition of a parking brake (or a side brake); a potentiometer 146 for detecting the position of the operation condition changeover valve 88; a valve opening degree sensor 148 for detecting the opening degree of a throttle valve (not shown) of the engine 2; an engine speed sensor 150 for detecting the speed of the engine 2; a drum speed sensor 152 for detecting the speed of the kick down drum 52 (FIG. 1); a gear speed sensor 154 for detecting the speed of the driven gear 64, to thereby detect the speed of the output shaft 50 and the speed of the vehicle; a pedal switch 156 which is turned on when the brake pedal of the vehicle (not shown) is depressed (or when the foot brake is operated); a low speed switch 158 which, when forward movement operating conditions hold, fixedly holds the speed change ratio in a lower speed range such as the second speed range; and a release switch 159 for releasing the parking mechanism 59 in priority irrespective of the operating conditions. Furthermore, the on-off signals of the vehicle key switches are applied to the electronic control device 112. The indicating switch 138 is made up of a conventional push switch incorporating a relay. The operating condition of the switch is restored to the initial operating condition when the switch is operated twice. The switch 138 has restoring means which operates to set the switch 138 to the OFF position, i.e., the forward movement indicating position in response to a signal from the electronic control device 112. The gear speed sensor 154 operates to detect when the speed of the driven gear 64 reaches 0 r.p.m. or approximately 0 r.p.m., thereby to detect a vehicle speed of 0 km/h; that is, the sensor 154 can detect when the vehicle stops. Thus, the gear speed sensor 154 serves as a vehicle stop sensor. The pressurized oil, which is discharged from the oil pump 74, is delivered through the oil path 160 to the vent valve 82, the operating condition change-over valve 88 and the reducing valve 86. The operating condition change-over valve 88 is employed in place of the manual valve in the conventional automatic speed change gear assembly, and comprises a spool 164 having a threaded groove 162 at the outer end portion, a servo motor 170 having a rotary shaft 168 fixedly secured to a disk 166 engaged with the threaded groove 162, the servo motor 170 being operated in response to an electrical signal from the electronic control device 112, and a potentiometer 146 coupled to the spool 164, to feed back an electrical signal corresponding to the position of the spool 164 to the electronic control device 112. The operating condition change-over valve 88 is controlled by the electronic control device 112 according to a control process described hereinafter. The change-over valve 88 has four positions; D, N, R and P. When the valve 88 is at the position D, the oil path 160 communicates with oil paths 172 and 174, to allow the speed change gear assembly 22 to provide the forward movement operating conditions in the first through fourth speed ranges according to the on-off combinations of the solenoid valves 108 and 110 as will be described later. When the valve 88 is at the position N, the oil path 160 communicates with the oil path 174 only, and the oil path 172 communicates with an oil discharging outlet 176, so that the speed change gear assembly 22 is placed in the neutral state. When the change-over switch is set to the position R, the oil path 160 communicates with oil paths 178 and 180, so that the speed change gear assembly 22 provides the reverse range. When the change-over switch 88 is set to the position P, all the oil paths to the change-over switch 88 are communicated with the oil discharging outlet 176 or 182, so that the speed change gear assembly 22 is placed substantially in the neutral state. Simultaneously when the signal for setting the operating condition change-over valve 88 to the position P is produced by the electronic control device 112, a signal is applied to the actuator 56 (FIG. 1) by the electronic control device 112. As a result, the actuator 56 is driven, so that the pawl 58 is engaged with the engaging teeth 54 to stop the rotation of the output shaft 50, whereby the vehicle is placed in the parked state. The vent valve 82 comprises a spool 188 having pressure receiving surfaces 184 and 186, and a spring 190. When the oil pressure in the oil path 160 acts on the pressure receiving surface 184 through the oil path 174, the oil pressure in the oil path 160 is adjusted to a constant value of 6 kg/cm 2 (hereinafter referred to as the line pressure) by the vent valve 82. When the oil pressure in the oil path 160 acts on the pressure receiving surface 186 through the oil path 178, the oil pressure in the oil path 160 is adjusted to 14.6 kg/cm 2 . The torque converter control valve 84 comprises a spool 192 and a spring 194. The control valve 84 operates to adjust the pressurized oil introduced through an oil path 196 from the vent valve 82, to 2.5 kg/cm 2 by balancing the oil pressure acting on the right pressure receiving surface of the spool 192 through a passage 198 formed in the spool 192 and the energization force of the spring 194, and to apply the pressurized oil thus adjusted to the torque converter 6 through an oil path 200. The oil which is discharged from the torque converter 6 is supplied through an oil cooler 202 to lubricate the various parts of the speed change gear assembly. The reducing valve 86 comprises a spool 204 and a spring 206. The reducing valve 86 operates to adjust the oil pressure from the oil path 160 to 2.4 kg/cm 2 by balancing the oil pressure due to the difference in area between the opposed pressure receiving surfaces 208 and 210 of the spool 204 and the energization force of the spring, and to apply the oil pressure thus adjusted to an oil path 212. The pressure-adjusted oil introduced to the oil path 212 is supplied through an orifice 214 to the N-R control valve 94, the oil pressure control valve 96 and an orifice 114 of the solenoid valve 106. The N-R control valve 94 comprises a spool valve 222 having pressure receiving surfaces 216, 218 and 220, and a spring 224. The control valve 94 operates to adjust the oil pressure of the oil path 226 to a predetermined value by balancing the oil pressure acting on the pressure receiving surface 216 and the composite force of an oil pressure due to the difference in area between the pressure receiving surfaces 218 and 220 and the energization force of the spring. The oil pressure control valve 96 has a spool valve 234 with pressure receiving surfaces 228, 230 and 232, and a spring 236. The control valve 96 operates to adjust the oil pressure of an oil path 238 to a predetermined value by balancing the oil pressure acting on the pressure receiving surface 228, and the composite force of an oil pressure due to the difference in area between the pressure receiving surfaces 230 and 232 and the energization force of the spring 236. The adjusted oil pressure, which is introduced to the oil path 226, is for controlling the low reverse brake 32 is providing the reverse range. The adjusted oil pressure, which is introduced to the oil path 238, is for controlling the front clutch 24, the rear clutch 26, the kick down brake 30 and the low reverse brake 32 when the vehicle is being run forwardly or stopped. The solenoid valve 106 is subjected to duty control with a constant frequency pulse current of 50 Hz, the pulse width of which is changed according to the operating conditions, by the electronic control device 112. That is, the pulse width is changed, to vary the ratio of opening time to closing time of the orifice, to thereby control the oil pressure in the oil path 212 downstream of the orifice 214, i.e., the oil pressure P 1 applied to the pressure receiving surface 216 of the N-R control valve 94 and to the pressure receiving surface 228 of the oil pressure control valve 96. For instance, in the case where the diameters of the orifices 214 and 114 are 0.8 mm and 1.4 mm, respectively, the oil pressure P 1 is adjusted in the range of about 0.3 to 2.1 kg/cm 2 . As the oil pressure P 1 is increased or decreased as described above, the adjusted oil pressure in the oil paths 226 and 238 is proportionally increased and decreased in the range of from about 0 kg/cm 2 to the supply oil pressure (which is the oil pressure in the oil path 180 or 172). The operation start timing and the operating period of the solenoid valve 106 described above are determined according to the electrical signals provided by the throttle valve opening degree sensor 148, the speed sensors 150, 152 and 154, a speed change detecting unit incorporated in the electronic control device 112 to detect the start of a speed change and an engaging timing detecting device comprising two speed sensors 152 and 154. The shift control valve 90 is controlled by combining the opening and closing operations of the solenoid valves 108 and 110. The shift control valve 90 has three spools 240, 242, 244 and two stoppers 246 and 248. The spool 240 has lands 250 and 252, an annular groove 254 and an oil path 258 which communicates the groove 254 with an oil chamber 256 on the left side of the land 250. The spool 242 has lands 260 and 262 different in diameter from each other, an annular groove 264, and abutting parts 266 and 268 which are abutted against the spools 240 and 244, respectively. The spool 244 has lands 270 and 272, an annular groove 274, and an oil path 278 communicating the groove 274 with an oil chamber 276 on the right side of the land 272. The stoppers 246 and 248, being held between the spools 240 and 242; are fixedly secured to the casing. The oil path 172 communicates with an oil path 280 through the annular groove 264 at all times. The oil path 280 communicates through an orifice 282 with the orifice 116, the left oil chamber 256 and the right oil chamber 276, and through an orifice 284 to an orifice 118 and an oil chamber 286 between the spools 240 and 242. The combination of the opening and closing operations of the solenoids 108 and 110 are related to the speed change ranges as shown in FIG. 4. The rear clutch control valve 92 has a spool 294 which includes a land 288, a land 290 smaller in diameter than the land 288, and an annular groove; a spool 306 which includes three lands 296, 298 and 300 equal in diameter to the land 290, and annular grooves 302 and 304; and a spring 308. When the oil pressure introduced into an oil chamber 310 (on the upper right side in FIG. 3A) to act on the pressure receiving surface of the land 288 becomes larger than the composite force of the oil pressure introduced to an oil chamber 312 (on the upper left side in FIG. 3A) to act on the pressure receiving surface of the land 300 and the energization force of the spring 308, the spools 294 and 306 are displaced to the right in FIGS. 3A and 3B. Upon displacement of the spools 294 and 306 to the right, the oil pressure is applied between the lands 290 and 294. Therefore, when the oil pressure is removed from the oil chamber 310, only the spool 294 is moved to the left. Thereafter, when the oil pressure applied to the left pressure receiving surface of the land 296 becomes smaller than the composite force of the pressure in the oil chamber 312 and the force of the spring 308, the spool 306 is moved leftwardly. The N-D control valve 98 has a spool 320 which includes lands 314 and 316 and an annular groove 318, and a spring 322. The spool 320 is displaced selectively between the left end position shown in FIG. 3B and the right end position according to the direction of the composite force of oil pressures applied to the pressure receiving surfaces 324, 326 and 328 of the spool 320 and the force of the spring 322. The first-second speed shift valve 100 comprises a spool 330 and a spring 332. The spool 330 is displaced between the left end position as shown in FIG. 3B and the right end position depending on whether or not the line pressure is applied to the left pressure receiving surface 334 of the spool. When the line pressure is supplied to act on the pressure receiving surface 334, the spool 330 is moved to the right end. When the line pressure is removed, the spool is moved to the left by spring 332. Similarly, the second-third and fourth-third speed shift valve 102 and the four-speed clutch control valve 104 have a spool 336 and a spring 340, and a spool 338 and a spring 342, respectively. Oil chambers 344 and 346 to which the line pressure is introduced are formed on the left sides of the spools 336 and 338, respectively, and oil chambers 348 and 350 are formed on the right sides thereof, respectively. These spools are displaced selectively to either the left end positions as shown in FIG. 3B or to the right end positions. Speed change control by selective engagement of the frictional engaging devices will now be described. When the electronic control device 112 drives the servo motor 170 according to a control process (described hereafter) until the spool 164 of the operating condition change-over valve 88 is set to the position D, the line pressure in the oil path 160, which has been adjusted to 6 kg/cm 2 , is conducted through the oil path 172 to the shift control valve 90. If, in this case, the solenoid valves 108 and 110 are electrically energized, the orifices 116 and 118 are open. Accordingly, the oil pressures in the oil chambers 256, 276 and 286 are not increased, the spool 242 is set to the left end position (as shown in FIGS. 3A and 3B) by the oil pressure which is provided according to the difference in pressure receiving surface area between the lands 260 and 262, and the line pressure is not conducted to other oil paths communicating with the shift control valve 90. The line pressure in the oil path 172 is led through the oil pressure control valve 96, the oil path 238, the N-D control valve 98, the oil path 352, the rear clutch control valve 92 and the oil path 354 to the oil pressure chamber of the rear clutch 26, and is led through the oil path 238, the first-second speed shift valve 100 and the oil path 356 to the oil pressure chamber of the low reverse brake 32. Thus, the first speed range is obtained by engagement of the rear clutch 26 and the low reverse brake. In this case, the operation of the oil pressure control valve 96 prevents an abrupt rise of the oil pressure in the oil pressure chamber of the rear clutch 26, to reduce the speed change shock. The operation of the N-D control valve 98 in this operation is as follows: When the oil pressure is led from the oil path 358 to the annular groove 318, the spool 320 is displaced to the right hand side in FIG. 3B against the elastic force of the spring 322 by the difference in area between the pressure receiving surfaces 326 and 328, as a result of which the communication of the oil path 238 with the oil path 352 is interrupted, and the oil path 172 communicates through an oil path 358 with an oil path 352 and with an oil path 360. Accordingly, as long as the operating condition change-over valve 88 is set to the position D, the oil path 172 communicates with the oil path 352 while bypassing the oil pressure control valve 96. Accordingly, the oil pressure applied to the oil pressure chamber of the rear clutch 26 is not affected by the oil pressure reduction control effected by the oil pressure control valve 96 at the time of speed change, and therefore the occurrence of difficulties such as speed change shock due to the slip of the rear clutch 26 in changing the speed range and the idling of the engine can be prevented. When, under the condition that the first speed range has been obtained, the acceleration pedal is depressed to increase the speed of the vehicle, the electronic control device 112 issues a second speed range instruction in response to signals from the throttle valve opening degree sensor 148 and the gear speed (vehicle speed) sensor 154, so that the solenoid valve 108 is deenergized, while the solenoid valve 110 is maintained energized. By this switching operation, the line pressure in the oil path 280 is led through the orifice 282 to the annular groove 254, the oil path 258, the oil chambers 256 and 276 and the annular groove 274. As a result, the spool 240 together with the spool 242 is moved to the right and is stopped in abutment against the stopper 246, whereupon the line pressure in the oil path 172 is conducted through the annular groove 264 to the oil path 362 and is applied to the pressure receiving surface 334 of the first-second speed shift valve 100 and the pressure receiving surface 346 of the four-speed clutch control valve. As a result, spools 330 and 338 of the valves 100 and 104 are moved to their right end positions. As a result, the line pressure in the oil path 238 is applied through the oil path 364 to the oil chamber 366 on the engagement side of the kick down brake 30 so as to move the rod 368 to the left against the elastic force of the spring 370, so that a brake band (not shown) is engaged with the kick down drum 52. At the same time, the oil pressure in the oil path 356 is discharged through the oil path 226 to release the engagement of the low reverse brake 32. Thus, the second speed range is obtained. In this operation, the oil pressure control valve 96 reduces, in response to the oil pressure control operation of the solenoid valve 106, the oil pressure in the oil path 238, i.e., the oil pressure which is applied to the oil chamber 366 of the kick-down brake 30, thereby to prevent speed change shock. When the application of current to the solenoids 108 and 110 is interrupted in order to obtain the third speed range in response to an instruction from the electronic control device, the line pressure is applied to the oil chamber 286 through the orifice, the spool 242 is moved right in FIG. 3 by the line pressure acting on the pressure receiving surface of the land 260 and is stopped in abutment against the spool 244, and the oil path 172 communicates with an oil path 372. The line pressure led to the oil path 372 acts on the pressure receiving surface 344 of the second-third speed and fourth-third speed shift valve 102, so that the spool 336 of the valve 102 is moved to the right and therefore the oil path 364 communicates through the orifice 374 with the oil path 376. The oil pressure led to the oil path 376 is supplied through the change-over valve 378 and the oil path 380 to the oil chamber 350 of the four-speed clutch control valve 104, and to release side oil chamber 382 of the kick-down brake 30 and through the change-over valve 384 to the front clutch 24. As the oil path 376 communicates with the oil chamber 382 of the kick-down brake 30 and the front clutch 24, engagement and disengagement of the two elements are carried out in an overlapping fashion. In shifting from the second speed range to the third speed range, the oil pressure control valve 96 is operated completely similarly as when shifting from the first speed range to the second speed range, so that the oil pressure in the oil path 238 is maintained at a low level for a short period of time. The oil path 364, which communicates through the first-second speed shift valve 100 to the oil path 238, is coupled to an orifice 374. During the operation of the oil pressure control valve 96, by the action of the orifice 374, the oil pressures in the release side oil chamber 382 of the kick-down brake 30 and the oil chamber of the front clutch 24 are maintained equal to one another, engagement of the front clutch 24 is carried out in parallel with the release of the brake 30, and thereafter when the operation of the oil pressure control valve 96 is stopped to raise the oil pressure to 6 kg/cm 2 the engagement of the front clutch 24 is accomplished. Thus, the third speed is obtained. In this case, the rotational speed of the input shaft 20 and the kick-down drum 52 approaches that of the output shaft, and finally the former becomes equal to the latter. The time instant when the two rotational speeds become equal, or immediately before the two rotational speed become equal, being regarded as the time instant when shifting has been accomplished, is detected by the speed sensor 152 and 154. Upon detection of this time instant, the operation of the oil pressure control valve 96, i.e., the operation of the solenoid valve 106 is stopped, to raise the supply pressure to the front clutch 24 to 6 kg/cm 2 . By this pressure increase, the oil pressure in the oil chamber 350 of the four-speed clutch control valve 104 is also increased, so that the spool 338 is displaced to the left end position in FIG. 3, the line pressure in the oil path 372 is applied through the oil path 386 to the four-speed clutch 28, and the latter 28 is placed in its engagement state. The oil path 386 communicates through a switching valve 378 and an oil path 380 to the oil chamber 350. Once the line pressure is supplied to the oil path 386, the spool 338 of the four-speed clutch control valve 104 is maintained at the left end position in FIG. 5 until the line pressure is removed from the oil path 386, to thereby prevent a difficulty wherein, in shifting from the third speed range to the fourth speed range or vice versa, the clutch 28 is released or slips whereby the shifting of the speed range becomes impossible or the neutral state is provided. In order to obtain the fourth speed range in response to an instruction from the electronic control device 112, the solenoid valve 108 is energized and the solenoid valve 110 is maintained deenergized. The oil pressures in the oil chambers 256 and 276 of the shift control valve 90 are thus decreased, so that the spool 244 together with the spool 242 is moved to the right end position in FIGS. 3A and 3B. As a result, the line pressure in the path 172 is led through the oil path 388 to the oil chamber 310 of the rear clutch control valve 92 and through the check valve 390 to the oil path 386. The spools 294 and 306 of the rear clutch control valve 92 are moved to the right end positions in FIG. 3B by the line pressure applied to the oil chamber 310, the oil path 388 is thus communicated with the oil path 392, the line pressure is supplied to the oil chamber 348 of the second-third speed and fourth-third speed shift valve 102, and the spool 336 of the valve 102 is moved to the left end position FIG. 3B. In this operation, the oil pressure in the oil chamber of the rear clutch 26 is discharged through the oil discharge outlet 394 of the rear clutch control valve 92, so that the rear clutch 26 is immediately released. On the other hand, the oil pressure in the oil chamber of the front clutch 24 and the oil pressure in the oil chamber 382 of the kick-down brake 30 are eliminated through the oil discharge outlet 396 of valve 102, so that the front clutch 24 is released and the kick-down brake 30 is engaged. Similarly as in the case of changing from the first speed range to the second speed range or from the second speed range to the third speed range, the oil pressure control valve 96 is operated to reduce the oil pressure in the oil path 238 for a short period of time during the speed changing operation, whereby the oil pressure acting on the oil chamber 266 of the kick-down brake 30 is decreased, so that the engagement thereof is carried out smoothly. Thereafter, when this oil pressure is increased to 6 kg/cm 2 , the engagement is achieved. Thus, the fourth speed range has been obtained. Now, downshifting will be described. The procedure of switching the operating paths of the oil pressure is opposite that in the case of the above-described upshifting sequence. In order to change the fourth speed range to the third speed range in response to an instruction from the electronic control device 112, the solenoid valves 108 and 110 are deenergized, as a result of which the line pressure in the oil path 388 is eliminated. In this case, the line pressure from oil path 352 is applied to the left pressure receiving surface of the land 296 of the rear clutch control valve 92, and similarly as in upshifting, the oil pressure control valve 96 is operated to reduce the oil pressure in the oil path 238. Therefore, only the spool 294 is moved to the left end position, so that the oil pressure in the oil path 392 is eliminated, and the spool 336 of the second-third speed and fourth-third speed shift valve 102 is moved to the right position. As a result, the oil pressure in the oil path 364 is smoothly applied through the orifice 374 and the oil path 376 to the oil chamber 382 of the kick-down brake 30 and the front clutch 24. When release of the kick-down brake 30 and engagement of the front clutch 24 are substantially achieved, the operation of the oil pressure control valve 96 is stopped and the oil pressure in the oil path 238 is increased. Therefore, the spool 304 of the rear clutch control valve 92 is moved to its left end position, so that the line pressure from the oil path 352 is supplied through the oil path 354 to the rear clutch 26 whereby the latter is engaged. Thus, the third speed range is obtained. The rear clutch 26 is thus engaged later than the front clutch 24. This is to reduce the speed change shock which is caused when the rear clutch, which is larger in torque capacity, is engaged first. When the third speed range is changed to the second speed range, the solenoid valve 108 is deenergized, while the solenoid valve 110 is energized, so that the oil pressure in the path 372 is eliminated. Therefore, the spool 338 of the four-speed clutch control valve 104 is moved to the right end position so that the oil pressure in the oil path 386 is discharged through the oil path 388, while the spool 336 of the shift valve 102 is moved to the left position so that the oil pressure in the oil path 376 is discharged. Thus, the four-speed clutch 28 and the front clutch 23 are released, while the kick-down brake 30 is engaged. Thus, the second speed range is obtained. In this speed change, similarly as in the above-described case, the oil pressure control valve 96 is operated, so that the kick-down brake 36 is smoothly engaged, thus preventing speed change shock. When the second speed range is downshifted to the first speed range, the solenoids 108 and 110 are energized. As a result, the spools 240 and 242 of the shift control valve 90 are moved to the left positions to remove the oil pressure from the oil path 362, while the spool of the four-speed clutch control valve 104 is moved to its left end position, to remove the oil pressure from the oil chamber 366 of the kick-down brake 30, whereby the latter is released and the low reverse brake 32 is engaged. Thus, the first speed range is obtained. When the electronic control device 112 outputs a signal to drive the servo motor 170 to thereby set the spool 164 of the operating condition change-over valve 88 to the position R, the oil path 160 communicates with the oil paths 178 and 180. The oil path 180 communicates through the N-R control valve 94, the oil path 226, the first-second speed shift valve 100 and the oil path 356 to the low reverse brake 32, while the oil path 178 communicates through the switching valve 384 to the front clutch 24. Therefore, the front clutch 24 and the low reverse brake 32 are engaged, to obtain the reverse range. In this case also, similarly to when changing the forward drive speed ranges, the solenoid valve 106 is operated for a short period of time, and the oil pressure applied to the low reverse brake 32 is maintained at a low level by the control of the N-R control valve 94, so that shock is prevented. The operating condition change-over valve 88 and the solenoid valves 108 and 110 are switched according to a control process shown in FIGS. 5A-5C by the electronic control device 112 in response to signals provided by the various switches and sensors. In FIG. 5A, reference character (1) designates a key switch "on" process. In process (1), the power switch is turned on, to perform a program start process (2). Thereafter, a release switch 159 is turned off, so as to prevent the occurrence of a problem wherein if the release switch 159 is in an "on" state, the parking mechanism 59 (indicated by reference character p59 in FIGS. 5A-5C) can be released immediately after the key switch "on" process (1). Thereafter, a data reading process (4) is carried out. Then, a release switch status determining process (5) for determining whether or not the release switch 159 is in an "on" state, is carried out. When the result of this process is "No," a process (6) for determining whether or not the speed of the engine is more than a predetermined value a 1 (for instance 600 r.p.m) is carried out. When the result of this determination is "No," then the reading process (4) is effected again. That is, processes (4), (5) and (6) are carried out again. When, during this operation, the operator turns on the release switch 159, a "Yes" is obtained as the result of process (5), so that process (7) for releasing the parking mechanism 59 is carried out, and a program ending process (8) is then performed. Thus, even when the key switch is turned off, the parking mechanism 59 remains released. This is useful when the vehicle must be towed, for instance, when the engine becomes out of order. While the processes (4) through (6) are repeatedly carried out as described above, the key switch may be set to the starter position, to operate the engine. When the speed of the engine reaches the predetermined value a 1 , a "Yes" is obtained as the result of the process (6). Therefore, a determining process (9) for determining whether or not the seat switch 142 is in an "on" state, a process (10) for determining whether or not the door switch 140 is turned on, and a process (11) for determining whether or not a side switch is turned on are carried out. The process (9) detects whether or not the operator sits on the seat; process (10) detects whether or not the door beside the operator's seat is closed; and process (11) determines whether or not the parking brake is operated. When a "No" is provided as a result of any one of the processes (9) through (11), a process (12) for determining whether the key switch is turned on or off is carried out. When "Yes" is provided as a result of this process, a process (13) is carried out in which the operating condition change-over switch 88 is moved to the position P or held there and the parking mechanism 59 is engaged, so that the vehicle is placed in parked state. Then, a data reading process (14) is carried out, and the determining process (9) through (14) (hereinafter referred to as process loop A) are again carried out. These processes are inserted to avoid any danger that the vehicle might start running immediately after the engine has been operated or that the vehicle may be caused to roll down a slope by inertia when the parking mechanism 59 is released. These processes allow for engine warm up before the vehicle starts running, and does not allow the auto to be run before the steps in which the operator sits down on the driver's seat and closes the door with the parking brake in operation have been achieved. If the engine is stopped while process loop A is being carried out, the starter should be operated to start the engine again. In this case, process loop A is repeated again. When the key switch is turned off during this period, "No" is obtained as a result of process (12). Therefore, a process (15) for placing the vehicle in a parking state (similarly as in process (13)) is carried out, and a program ending process (16) is effected. When the results of the determining processes (9), (10) and (11) are all "Yes," a process (17) for releasing the parking mechanism 59 and a process (18) for reading data are carried out, and a release switch status determining process (19) is carried out. When a "No" is provided as a result of process (19), i.e., when the release switch 159 is in the "off" state, process (20) for determining whether or not the vehicle is stopped, i.e., whether or not the vehicle speed is 0 km/h, is carried out. When process (20) is carried out immediately after processes (9), (10), (11), (17) and (18), "Yes" is provided as a result of process (20) because the vehicle speed is generally 0 km/h. In this case, the parking brake is normally in operation, and therefore "Yes" is obtained as a result of the side switch on status determining process (21). Accordingly, a process (22) for determining whether or not the speed of the engine is 0 r.p.m is performed. When the result of this determination is "No," a process (23) for setting the operating condition change-over valve 88 to the position N is conducted, and the data reading process (18) is effected. As long as the parking brake is in operation, the processes (18) through (23) (hereinafter referred to as process loop B) are repeated. When, during process loop B, the operator releases the parking brake and "No" is provided as a result of process (21), a seat switch "on" determining process (24) is conducted. When "Yes" is provided as a result of process (24), process (25) for determining whether or not the speed of the engine is 0 r.p.m. is carried out. When the result of this process (22) yields "No," an indicating switch status determining process (26) for determining whether the indicating switch 138 indicates forward movement or reverse movement, i.e., whether the indicating switch 138 is turned off or turned on, is performed. When "No" is provided as a result of process (26), i.e., when forward movement is indicated, a process (27) for setting the operating condition change-over switch 88 to the position D is performed and a process (28) is carried out to determine whether or not a low speed switch 158 for maintaining the change gear ratio in the low range is turned on. When "No" is provided as a result of this process (28), a process (29) for instructing a speed change is performed. In the speed change instructing process (29), the electronic control device 112 applies signals to the solenoid valves 108 and 110 to provide a forward speed range which is suitable for the operating conditions of the vehicle at that time. After process (29) has been achieved, a process (30) for determining whether or not the throttle valve is fully closed, a process (31) for determining whether or not acceleration of the vehicle is larger than a predetermined value C 1 , and a process (32) for determining whether or not the foot brake is in operation, i.e., whether or not the pedal switch 156 is turned on, are carried out. When "No" is provided as a result of any one of the processes (30), (31) and (32), the data reading process (18) is again effected. When the vehicle is run ordinarily and accelerated with the throttle valve open, when the vehicle is running down a slope with the throttle being fully closed and with the acceleration being maintained smaller than the predetermined value C 1 , and where the vehicle is being decelerated, processes (18), (19), (20), (25), (26), (27), (28), (29), (30), (31) and (32) (hereinafter referred to as process loop C) are repeatedly performed. When, during process loop C, the vehicle is stopped after being decelerated, i.e., when the vehicle speed becomes 0 km/h, a "Yes" is provided as the result of process (20), and a process loop C' in which the processes (21) and (24) are inserted between processes (20) and (25) of process loop C is effected. When, under this condition (i.e., where the vehicle is stopped), the parking brake is operated, the side switch 144 is thus turned on, and a "Yes" is provided as the result of process (21). Thus, process loop B is again effected, and the speed change gear assembly 22 is set in the neutral state. If the vehicle is stopped by the foot brake in process loop C', and the operator then leaves the vehicle without applying the parking brake, a "No" is provided as the result of the seat switch actuation determining process (24). Therefore, after completing determining process (12), process (13) for parking the vehicle is performed, to prevent the vehicle from being run without an operator. In this case, processes (12), (13), (14), (9) and (12) are effected. Therefore, even if the operator then sits on the driver's seat and a "Yes" is provided as the output of process (9), the operation sequence is not shifted over to process loops C or C' unless the parking brake is operated. When, during forward running with process loop C repeatedly carried out, the operator or a passenger turns on the indicating switch 138 (indicating reverse) by accident, a "Yes" is provided as the result of determining process 26. However, since "No" is provided as the result of process (33) (determining whether or not the vehicle speed is 0 km/h), process (23) for setting the operating condition change-over valve 88 to the N position is carried out and the speed change gear assembly 22 is placed in neutral. The processes (18), (19), (20), (25), (26), (33) and (23) are carried out until the indicating switch 138 is again switched to indicate forward driving or until the vehicle speed is reduced to 0 km/h. When the indicating switch 138 is turned off again, the above-described process loop C is again carried out and the forward speed range is reinstated. Accordingly, if the operator turns on the indicating switch 138 by mistake during forward driving of the vehicle, no reverse speed range is provided. This greatly contributes to the operational security of the vehicle. In the condition that the process loop C consisting of processes (18), (19), (20), (21), (24), (25), (26), (27), (28), (29), (30), (31) and (32) is being repeated, when the operator operates the foot brake because the vehicle is running down a slope with the throttle valve fully closed and with the vehicle's acceleration speed larger than the predetermined value C 1 , then the results of the determining processes (30), (31) and (32) are all "Yes." Thus, after a data reading process (34) is completed, a low speed range specifying process (35) for obtaining a lower speed range is carried out. In process (35), the electronic control device 112 applies signals to the solenoid valves 108 and 110 so that a lower speed range such as the third or second speed range is provided according to the speed or acceleration of the vehicle at that time. Thereafter, a process (36) for determining whether or not the engine speed is 0 r.p.m., a process (37) for determining whether or not the vehicle speed is 0 km/h, a process (38) for determining whether or not the low speed switch 158 is turned on (in this case, the low speed switch 148 is in the "off" state and "No" is provided as a result of this process because the process is effected after the process loop C), and a process (39) for determining whether or not the throttle valve is fully closed are carried out. As long as the engine is maintained running, the vehicle continues traveling and the throttle valve is fully closed, the processes (34), (35), (36), (37), (38) and (39) (hereinafter referred to as process loop E) are performed. Thus, even after the foot brake is released, the above provided second or third (downshifted) speed range is maintained, so that an engine braking action is provided. When, during the execution of process loop E, the vehicle is stopped or the throttle valve is again opened, a "Yes" is provided as the result of process (37) or a "No" is provided as the result of process (39), so that process loop C or C' is then effected. An input element to the electronic control device 112, namely, a cancel switch 400 for preventing the downshift to a lower range may be provided as indicated by the broken line in FIG. 3. In addition to the processes (30), (31) and (32) in this case, a cancel switch status determining process (40) for determining whether or not the cancel switch 400 is turned on, as indicated by the broken line in FIG. 5, is carried out in shifting to process loop E from process loop C. When "Yes" is provided as a result of the process (40) (i.e., when the cancel switch 400 is turned on), process loop E is not effected, and accordingly the downshift is not executed. During the execution of process loop C or C', if the operator turns on the low speed switch 158, a "Yes" is provided from determining process (28). Therefore, the data reading process (34) and the low speed range specifying process (35) are performed. In addition, determining processes (36) through (38) are carried out. The processes (34) through (38) (hereinafter referred to as process loop E') are then repeatedly carried out, in normal running, until the low speed switch 158 is turned off. Accordingly, the electronic control device 112 applies to the solenoid valve 108 and 110 signals for providing a low speed range in response to an electrical signal from the low speed range in response to an electrical signal from the low speed switch 158, or signals for achieving an automatic shift between the low speed range and a lower speed range. Accordingly, the engine braking action can be achieved as required by the operator, or quick acceleration can be conducted while holding the low speed range. In the process loop E', when the vehicle is stopped, a "Yes" is provided from process (37). Therefore, the data reading process (18) is effected, and processes (19) through (21) and (24) through (28) are performed. "Yes" is still provided as the result of process (28), so process loop E' is again carried out. Therefore, when the vehicle is stopped, processes (37), (18), (19), (20), (21), (24), (25), (26), (27), (28), (34), (35), (36) and (37) are repeated. Thereafter, when the vehicle is running, process loop E' is repeated. In addition to the indicating switching, a low speed switch 158 may be provided so that the second speed range, which is extensively employed as a low speed range, may be fixedly provided according to an electrical signal from this low speed switch. Alternatively, in the case of speed change gearing which, as in the present embodiment, includes four forward speed ranges, three low speed switches 158 may be provided in such a manner as to produce a signal for fixedly providing the first speed range, a signal for achieving automatic shifting between the first and second speed ranges, and a signal for achieving automatic shifting between the first, second and third speed ranges, respectively. When the operator turns on the indicating switch 138 (i.e., when reverse travel is instructed) when the vehicle is stopped and process loop C' is being repeated, the results of the determining processes (26) and (33) are "Yes." Therefore, a data reading process (41), a process for setting the operating condition change-over valve 88 to the position R, and a process (43) for sounding a buzzer or the like (optional) are performed. In succession, a process (44) for determining whether or not the engine is stopped, and an indicating switch status determining process (45) are conducted. As long as indicating switch 138 is in the "on" state, processes (41), (42), (43), (44) and (45) (hereinafter referred to as process loop F) are repeated, and the reverse range is effected in the speed change gearing assembly. If the indicating switch 138 is turned off when the vehicle is moving backwardly while process loop F is being repeated, the reverse range is maintained, because a "No" is provided from process (46) in determining whether or not the vehicle speed is 0 km/h, and the data reading process (41) is again effected. When the vehicle speed reaches 0 km/h, an alarm stopping process (47) is performed. The process loop C' is then effected, the operating condition change-over valve 88 is set to the position D, and the forward speed range is obtained. Accordingly, even if the indicating switch 138 is turned off by mistake while the vehicle is running in reverse, the reverse range is maintained. Accordingly, the danger that a forward speed range may be provided before the vehicle is stopped is eliminated, which again contributes to the operational security of the vehicle. When, in process loop F, a "Yes" is provided from the engine stop determining process (44), or if the engine is stopped, the alarm stopping process (47) is effected, and after processes (18), (19) and (20), or processes (18), (19), (20), (21) and (24) have been achieved, a "Yes" is provided at process (25). Therefore, a process (48) for automatically turning off the indicating switch 138 is carried out, and a process (49) for setting the operating condition change-over switch 88 to N position is performed, to provide the conditions for again starting the engine. Also in the other process loops B, C, C', E and E', when the results of engine stop determining processes (22), (25) and (36) thereof are "Yes," the above-described processes (48) and (49) are effected. After processes (48) and (49) are achieved, a data reading process (50) is carried out, and a process (51) for determining whether the keys switch is turned on or turned off, i.e., for determining whether or not engine stop was caused by the operator turning off the key switch, is performed. When a "Yes" output is provided from determining process (51), a process (52) for determining whether or not the speed of the engine is larger than a predetermined value a 1 (600 r.p.m. for instance) is conducted; and when not, the data reading process (50) is again effected, and processes (50), (51) and (52) are repeated. During this period, the operator may start the engine again with the key switch set to the starter position. When the speed of the engine is larger than the predetermined value a 1 for a predetermined period of time (two seconds for instance), a "Yes" is provided at process (52), and therefore a process (53) for determining whether or not the next speed of the engine is smaller than a predetermined value a 2 (1,500 to 2,500 r.p.m. for instance) is conducted. When a "No" is provided at process (53), the data reading process (50) is again effected; and when a "Yes" is provided, a process (54) for setting the operating condition change-over valve 88 to the position D is performed, and thereafter process loop C or C' is carried out. When the result of the key switch on status determining process (51) is "No," a process (55) for determining whether or not the vehicle speed is 0 km/h is carried out. After it is confirmed that the vehicle is stopped, the process (15) for parking the vehicle is carried out, and the program ending process (16) is performed. In the case where the engine stopped while the vehicle was running rearwardly, it is dangerous to allow the vehicle to start backward movement immediately after the engine is again started. In order to avoid this danger, the indicating switch status determining process (48) and the process (54) for setting the operating condition change-over valve 88 to the position D are carried out. That is, by the processes (48) and (50), the vehicle is shifted into a forward range initially, for security. When the engine stops while the vehicle is moving forwardly, and the operator starts the engine again while the vehicle is still moving due to inertia, the engine may stop again if a speed change is achieved while the rotation of the engine is unstable; and so-called "speed change shock" is caused if speed change is achieved while the speed of the engine is high. In order to avoid these difficulties and thereby achieve smooth speed changing, the engine speed determining processes (52) and (53) are performed. The foregoing embodiment of the invention, operating according to the above-described control process and using an oil pressure control system and an electronic control device, has the following merits: (a) The parking mechanism 59 is not released before the engine is operated, the operator sits on the driver's seat and closes the door and the parking brake is operated. Accordingly, the vehicle may not be run unintentionally immediately after the engine is started, and the parking mechanism 59 will not be released to allow the vehicle to run by inertia. Thus, the invention provides high security. (b) when it is confirmed that the key switch is turned off and the vehicle is stopped, or when it is confirmed that the operator has left the vehicle without operating the parking brake, the parking mechanism 59 is engaged. Accordingly, the vehicle cannot be caused to run without the operator. (c) The ordinary running conditions, including the application of the engine brake when the vehicle runs down a slope, can be obtained merely by the on-off operation of the indicating switch 138 controlled by the operator. Accordingly, with the present automatic speed change gear assembly, unlike the conventional one, an intricate lever operation and assembly is eliminated, the driving operation is simplified, operability and security are improved, and erroneous operations rarely occur. Even when, while the vehicle is running, the indicating switch 138 is erroneously operated to indicate the opposite direction, the speed range then specified by the indicating switch 138 is not obtained, i.e., a neutral state is achieved or the speed range in the direction of vehicle running is obtained. Thus, the invention prevents damage to the vehicle. (d) When the engine stops while the vehicle is running or is stopped, the operating condition change-over valve 88 is set to the N position if the key switch is not in the "off" state, and accordingly the gearing assembly is placed in the neutral state so that the engine can be started again. Accordingly, the speed change gear assembly is free from the troublesome operation of the conventional speed change gear assembly in which the operating lever must be set to the N or P position before the engine can be again started. That is, with the speed change gear assembly of the invention, immediately after the engine stops, the engine can be again started with the key switch set to the starter position. Thus, operation is considerably simplified. (e) When the parking brake is operated with the vehicle stopped with the engine operating, the operating condition change-over valve 88 is set to the N position and the speed change gearing assembly is placed in neutral. Therefore, when it is required to stop the vehicle in response to a traffic signal, the vehicle can be stopped by merely operating the side brake lever, which is effective in preventing creep and can prevent application of an excessive load to the torque converter or the like. (f) When the key switch is turned off, i.e., when the engine is stopped with the vehicle stopped, the parking mechanism 59 is automatically engaged. Accordingly, with the assembly of the invention, unlike the conventional apparatus, the danger of the operator leaving the vehicle without engaging the parking mechanism and the vehicle running unintentionally can be positively eliminated. As is apparent from the above description, according to the invention, in the present automatic speed change gearing assembly, the means for applying instructions from the operator to the automatic speed change assembly consists solely of the indicating switch for indicating the direction of advancement. Thus, the provision of a conventional operating lever for switching oil pressure circuits is unnecessary. Accordingly, a larger and unimpaired space can be provided for the operator, and vehicle operability and security are improved, while erroneous operations can be prevented. Furthermore, as the lever operation is eliminated, the driving operation is simplified by as much. In addition, in parking the vehicle, with the present automatic speed change gear assembly, vehicle stoppage and key switch turn-off are detected, whereby the parking mechanism is automatically operated and the vehicle is maintained parked. Accordingly, troublesome operations for operating the parking mechanism are eliminated, and the danger that the vehicle may run unintentionally because the operator has forgotten to perform some operation is prevented. That is, operability and security are improved. In the above-described embodiment, a low speed switch 158 which provides a signal for fixedly holding a particular low speed range irrespective of the operating conditions is employed. However, the switch 158 may be eliminated as the case may be, because it is required only when the operator wants to determine the speed range independently. In this case, the determining processes (28) and (38) can be eliminated from the flow chart of FIG. 5, when is then modified such that the process (29) is effected after process (27), and such that when the result of process (37) is "No," the process (39) is effected. In the above-described embodiment, the advancement direction indicating switch 138 is a conventional push-lock type switch. Instead of this switch, a restoration type normally open switch may be employed which has two closed positions to which the armature is tripped by the operating section, and one open position between the two closed positions. When the armature is tripped to one closed position, forward movement is indicated; and when it is tripped to the other closed position, backward movement is indicated. When the operating section is operated, the armature is tripped to one of the two closed position, but normally the armature is set at the open position. In this case, when forward movement or backward movement is indicated by the normally open switch, this indication is held in the electronic control device 112 until movement in the opposite direction is indicated. In the determining process (11) in FIG. 5, instead of determining whether or not the side switch 144 is turned on, a determination as to whether or not indication is provided by the normally open switch (or whether the normally open switch indicates forward movement or backward movement, discrimination between forward and backward movements not being necessary) is employed. In determining processes (26) and (45), it is determined whether the indication of the normally open switch (which is held in the electronic control device 112) is for backward movement or forward movement. In succession, in process (48), the indication of the normally open switch, which is held in the electronic control deivce 112, is forcibly set for forward movement. In the above-described embodiment, the technical concept of the invention is applied to an automatic speed change gear assembly in which the forward, backward and neutral operating conditions can be obtained by merely operating the indicating switch 138; however, it should be noted that the automatic parking mechanism of the invention is not limited thereto or thereby. That is, the technical concept of the invention may be applicable to an automatic speed change gear assembly in which the abovedescribed operating conditions are obtained by operating an operating lever similarly as in the prior art. In this case, the parking mechanism 59 with the acutator 56 is provided, the position P of the operating lever is eliminated, and only the parking mechanism 59 is automatically operated by the control device.

A fully automatic transmission for use in automobiles requires the use of only a single manually operable switch in the vehicle compartment for selecting the travelling direction. All other functions of the transmission, namely the selection of the speed range, including neutral and the application of a parking mechanism are performed by a hydraulc control system in combination with an electronic control system operating according to a predetermined process.

[0001] This application claims priority to provisional application Ser. No. 61/669,240, filed on Jul. 9, 2012, the contests of which are incorporated herein by reference in their entirety. BACKGROUND OF THE INVENTION [0002] This invention relates to nanostructure arrays for multifunctional surfaces that are mechanically robust and provide superior optical and wetting properties such as antireflection and superhydrophobicity/hydrophilicity. The invention also relates to the fabrication process for the nanostructure arrays. [0003] Nanostructured surfaces have been widely studies for their superior optical and wetting properties such as antireflection and superhydrophobicity/hydrophilicity [1, 2, 3]. Due to subwavelength feature size, nanostructured surfaces behave as an effective medium with gradually varying index of refraction. Such a surface can be used to suppress Fresnel reflection at material interfaces, thereby acting as an anti-reflection surface and allowing broadband light to pass through without reflection loses [4]. In addition, both hierarchical roughness of these structures and an intrinsic chemical property of the surfaces can induce artificial superhydrophobicity or superhydrophilicity, which can be applied as self-cleaning or anti-fogging surfaces, respectively [5, 6]. [0004] Recently, high aspect ratio (approximately 5) silica nanocone structures, demonstrating structural superhydrophilicity or, in combination with a suitable chemical coating, robust superhydrophobicity, and enhanced transmissivity, have been successfully fabricated directly on a fused silica substrate with few defects and a large pattern area [7]. By using interference lithography and multiple shrinking mask etching, the desired aspect ratio nanocone structures were created for optimizing the multifunctionality of the textured surface [8]. Although the prior art structures provided notable performance, the mechanical stability and properties of these prior art structures are not suitable to sustain mechanical impact (including finger touch) because of the high aspect ratio and isolated nature of the structures. [0005] It is an object of the present invention to provide an alternative type of nanostructure for multifunctional surfaces (anti-reflectivity, superhydrophobicity, superhydrophilicity and superoleophobicity) with high mechanical strength and high optical and wetting performance compared to existing nanostructures [7]. Another object of the invention is a fabrication process for the structures disclosed herein that is simple, and cost effective for manufacturing the multifunctional surfaces. SUMMARY OF THE INVENTION [0006] In a first aspect, the invention is method of fabricating nanotextured structures, including making a master mold having an array of tapered structures to be replicated. The master mold is pressed into a curable polymer supported on a substrate and the polymer is cured. Thereafter, the mold is detached from the cured polymer to form the nanotextured structure. In a preferred embodiment, the master mold is fabricated using lithography and etching. The master mold may be made of silicon or fused silica glass. The curable polymer may be heat curable or photo curable. It is preferred that an anti-adhesion layer be deposited on the master mold to prevent sticking. It is also preferred that the pressing step utilize vacuum assisted wetting. [0007] The master mold may be fabricated with positive resist or negative resist. The profile, height and period of the tapered structures to be replicated are selected to control surface characteristics. In yet another embodiment, re-entrant textured surfaces are formed on the replicated, nanotextured structure. The re-entrant textured surfaces may be made by pressing a hot plate into the replicated structure, or by pressing a rotating plate into the replicated structure. Conditions are optimized for perfect re-entrant structures. [0008] In another aspect, the invention is a nanotextured surface structure comprising periodic tapered nanohole or inverted nanocone arrays. The surface may exhibit superhydrophilicity and when the surface includes a selected coating, the structure exhibits hydrophobicity. The surfaces of the invention may exhibit enhanced optical transmission and low reflectance. BRIEF DESCRIPTION OF THE DRAWING [0009] FIG. 1 is a schematic illustration of the process disclosed herein according to an embodiment of the invention. [0010] FIG. 2 is a scanning electron micrograph of an array of replicated inverted nanocone structures. [0011] FIG. 3 is a schematic illustration of imprinted inverted concave nanocone structures with a master mold fabricated with positive resist. [0012] FIG. 4 is a schematic illustration of imprinted inverted convex nanocone structures with a master mold fabricated with negative resist. [0013] FIG. 5 is a schematic illustration of the process of making re-entrant nanostructures for transparent superoleophobic surfaces. [0014] FIG. 6 is a schematic illustration of a window fabricated with the process according to an embodiment of the invention. DESCRIPTION OF THE PREFERRED EMBODIMENT [0015] The nanostructures disclosed herein are fabricated using a novel and advanced replication technology. With reference first to FIG. 1 , a master mold 10 is fabricated using any lithography technique and subsequent etching steps with any material, such as silicon or fused silica glass. Thereafter, the master mold is used to fabricate inverted nanocone arrays. The master mold 10 contents and is pressurized onto resist 12 that may be a heat or photo curable polymer on a glass substrate 14 . After curing the polymer resist 12 either with heat or ultraviolet (UV) light, the master mold 10 is detached from the replicated polymer surface 16 . Thus, nanocone arrays on the master mold 10 are inversely replicated into the polymer 12 , thus the surface is textured with nanohole arrays 18 as shown in FIG. 2 . During the imprint process, it is preferred that an anti-adhesion layer be coated on the master mold 10 surface to prevent sticking problems. Further, in order to have a perfect pattern transfer, and without any bubbles in the polymer, vacuum assisted wetting is used when the nanocone arrays are pressed into the liquid polymer 12 , which is also one of the important processes in this invention. [0016] Depending on the geometry of the master mold 10 , different inverted nanocone structures (or nanohole structures) can be fabricated. FIG. 3 shows a schematic illustration of an imprinted inverted concave nanocone structure with a master mold 10 fabricated with positive resist. When the mold is fabricated with positive resist [7], the convex nanocone structures are realized for the master mold. FIG. 4 shows imprinted inverted convex nanocone structures with the master mold 10 fabricated with negative resist. The profile, height and period of the imprinted pattern can be easily controlled, optimized for superhydrophobicity, superhydrophilicity, anti-reflectivity and mechanical robustness. For example, the structures in FIG. 3 tend to be mechanically more robust than those in FIG. 4 , but those in FIG. 4 may have better transparency and superhydrophobicity. [0017] Deng et al. recently fabricated transparent and superoleophobic surfaces. However, the surfaces consisting of accidental nanoparticles limited special-phase coherence, resulting in non-perfect periodic/quasi random structures that scatter light and make themselves a lot less transparent than glass [9]. [0018] We also disclose novel and simple fabrication methods for surfaces with enhanced transparency and oil-repellency or superoleophobicity (for anti-fingerprint surfaces) by creating re-entrant textured surfaces as shown in FIG. 5 . First, an imprinted surface with inverted nanocone arrays is fabricated as discussed above. The imprinted pattern is contact or pressurized by a hot plate 20 and the re-entrant structures are created by local reflow of the material of the imprinted pattern itself. An alternate embodiment is that the imprinted pattern is contacted and pressurized by a rotating plate 22 thereby creating the re-entrant structures 24 shown in FIG. 5 . Temperature and pressure of the hot plate, and pressure, friction and speed of the rotating plate can be optimized for perfect re-entrant structures. [0019] The aspect ratio and shape of the nanohole structures disclosed herein can be optimized to achieve better wetting or optical functions using the fabrication method disclosed herein. By texturing subwavelength nanoholes on both sides of glass and modifying their surface energies, it is possible to combine high-pressure robustness of superhydrophobicity and near-perfect transparency (or anti-reflection property). In addition to the synergetic effect, the nanohole surface can show a macroscopic anti-fogging function for practical applications, including transparent windshields and goggles that are self-cleaning outside. FIG. 6 shows potential applications of glass treated on both sides. This window 26 with the nanohole or inverted nanocone structure will enhance transmission, and at the same time, make the surfaces both superhydrophobic and superhydrophilic such as self-cleaning (outside) and anti-fogging (inside) glass. In addition, transparent anti-fingerprint and scratch resistant surfaces can be fabricated with these structures. [0020] The disclosed fabrication process, according to the invention, can also be used to produce a protective glass for digital cameras. In this case, nearly 100% of the incident light with a wide angle can be collected without any loss, so that pictures with better quality can be taken, even at night. Using such high transmission protective glass will also eliminate interference effects. Dust or contamination on the surface can be easily removed with water due to the superhydrophobic surface. Those of ordinary skill in the art will recognize that the application is not merely for digital cameras, but also for solar cells, because the efficiency of a solar cell is highly influenced by dust contamination and Fresnel reflection losses. In humid conditions, the technology disclosed herein can be used for anti-fogging glass and good images may be taken without any bubbles on the lens. [0021] The nanohole or inverted nanocone structures disclosed herein have stronger mechanical properties than tapered nanocone structures known in the prior art since nanohole structures are expected to withstand larger external forces, since they are now connected to each other and support each other in a two-dimensional network. Surfaces according to the invention with these nanohole structure also become scratch resistant. Once use would be as a scratch resistant surface for touch screen electronic devices. Anti-fingerprint surfaces can be realized by superoleophobicity resulting from its geometry combined with proper chemical coating on the nanohole structures with re-entrant textured surfaces in FIG. 5 . [0022] The frabrication method disclosed herein is compatible and can be adapted to all conventional 2D lithography techniques. The mold 10 with nanocone structures can be fabricated with any lithographic process. Moreover, self-assembly approaches such as colloidal lithography or block copolymer can be used to pattern the nanocone structures. [0023] The present invention creates nanohole or inverted nanocone structures with a tapered profile at the same time. These tapered nanohole or inverted nanocone structures with re-entrant geometry on the tips are suitable for enhancing transmission (anti-reflectivity) and wetting properties (superhydrophobicity, superhydrophilicity, and superoleophobicity). [0024] The multifunctional nanotextured surfaces disclosed herein consist of perfect periodic tapered nanohole arrays or inverted nanocone arrays. The mechanical robustness of the structures disclosed herein is higher than for conventional moth eye-like structures, arrays of isolated tapered nanostructures. In the present invention, the inverted nanocone structures are structurally supported by adjacent nanostructures, as was shown in conjunction with FIG. 2 . For applications such as touch screen panels or car windows, the proposed multifunctional nanotextured surfaces which have high mechanical robustness are more desirable in terms of durability and stability, since the nanostructures will be exposed to a lot of mechanical impacts. [0025] The multifunctionality of the surfaces disclosed herein can be used for a wide range of applications, including self-cleaning and omnidirectional optical elements, as well as anti-fogging optical lenses for microscopes in humid biological environments, and semiconductor lithography equipment. [0026] The numbers in square brackets refer to the references listed herein. [0027] It is recognized that modifications and variations of the present invention will occur to those of ordinary skill in the art, and it is intended that all such modifications and variations be included within the scope of the appended claims. REFERENCES [0028] [1] Y. F. Huang, et al. Nature Nanotechnology, 2, pp. 770-774, (2007) [0029] [2] C.-T. Wu et al., Chemistry of Materials, 22, pp. 6583-6589, (2010) [0030] [3] C. Suet et al., Applied Surface Science, 253, pp. 2633-2636, (2006) [0031] [4] K. Choi et al., Advanced Materials, 22, pp. 3713-3718, (2010) [0032] [5] X. Gao et al., Advanced Materials 19, pp. 2213-2217, (2007) [0033] [6] E. Martines et al., Nano Letters, 5, pp. 2097-2103, (2005) [0034] [7] K. C. Park et al., ACS Nano, 6, pp. 3789-3799 (2012) [0035] [8] H. J. Choi et al., U.S. patent application Ser. No. 61/477,792, April 2011 (MIT Case No. 14861) [0036] [9] Deng et al., Science, 335, pp. 67-70 (2012)

Inverted Nanocone Structures and Its Fabrication Process. The method of fabricating nanotextured structures includes making a master mold having an array of tapered structures to be replicated. The master mold is pressed into a curable polymer supported on a substrate and the polymer is cured. Thereafter, the mold is detached from the cured polymer to form the nanotextured structure.

BACKGROUND OF THE INVENTION 1. Field of the Invention The present invention generally relates to a method for producing foil electrode for use in aluminum electrolytic capacitor. Particularly, it is concerned with a method for etching an aluminum foil which has an excellent characteristic as an electrode. 2. Description of the Prior Art As is well known, the aluminum foil employed as an electrode for electrolytic capacitor must have a coarse surface of an extended effective surface area for the purpose of miniaturizing the size of, minimizing manufacturing cost of and/or improving the electrical characteristics of the products. Treatments of the aluminum foil for making it to have the coarse surface of the extended effective surface area are usually referred to as "etching process" and have customarily been performed in an aqueous solution by chemically or electrochemically. In this field of etching process, there has hitherto been made a variety of studies and a number of proposals for attaining great industrial advantages. As a result of the detailed analysis of these studies, the following four principal factors, which greatly influence the performance of the etching process, i.e., extending rate of the effective surface area of the etched foil, are pointed out, and individual technical know-hows are established for each of the four factors: 1. Aluminum foil (purity of the metal; species, amount and mode of dispersion of impurities in the metal; orientation, size and shape of the crystal grain of the metal; and oxide film covering the surface, state and thickness), 2. Etching solution (species and concentrations of principal ingredient and additives; and temperature), 3. Power supply (direct current, alternating current, direct current superimposed with alternating current, waveform; and current density), and 4. Auxiliary treatment (pretreatments, intermediate treatments and combinations of these treatments). Although it has been known that all of these factors can separately and independently influence the performance of the etching process and the configuration of the etched surface of the aluminum foil, it is confirmed that what is highly important for the performance and the configuration is the complexities, i.e., the interactions or mutual dependencies between and among these individual factors. In the other words, one of these factors in a given state may adversely interfere or synergetically cooperate with the other factor to create a combined influence on the performance of the configuration. SUMMARY OF THE INVENTION It is therefore the principal object of the present invention to provide a method for producing foil electrode for use in an aluminum electrolytic capacitor which has an excellent electrical characteristics. It is another object of the present invention to provide an etching process on an aluminum foil which gives the foil a large effective area while maintaining a good mechanical strength of the foil. In order to attain the stated objects, the present inventors have explored the previously described various factors involved in the formation, distribution, configuration and growth of the etching pits by observing them with an electron microscope and the like, keeping their minds to the points: 1. To raise the density of the etching pits as high as possible and to distribute them over the entire surface as even as possible, 2. To make the diameter of each etching pit as small as possible, say, for instance, 0.2 μm, and as uniform as possible, and 3. To form the etching pits to an extent that they do not grow too deep to penetrate the foil completely in order to maintain the mechanical strength of the foil. As a result of the exploration, it is finally confirmed that a maximum performance of the etching process can be expected under conditions described later in connection with the preferred embodiment. DESCRIPTION OF THE DRAWINGS In the drawings; FIG. 1 is a schematic diagram showing an arrangement for the electrodes in an etching tub, and FIGS. 2-5, are graphic representations of the influences of the current density and the bath temperature on the electrostatic capacitance, those of the species and concentrations of the additives on the electrostatic capacitance, those of the concentrations of the hydrogen chloride and the bath temperatures on the electrostatic capacitance, and those of the frequency of the power supply at the etching process and formation voltages on the electrostatic capacitance, respectively. DESCRIPTION OF THE PREFERRED EMBODIMENT According to the present invention, there is provided an improvement in or relating to the method of enlarging effective surface area of aluminum foil for use in electrolytic capacitor by electromechanically etching said aluminum foil in an aqueous etching solution including hydrogen chloride and at least one additives capable of anodizing, that is, forming porous oxide film on, the surface of said aluminum foil with power supplied in alternating current, which comprises; the concentration of the hydrogen chloride in said solution is set at a range between 0.5 N and 3.5 N, the total concentration of said additives in said solution is set at a range between 0.001 mole/liter and 0.4 mole/liter, and said alternating current is being supplied at a frequency ranging from 5 Hz to 120 Hz and in a current density ranging from 0.1 A/cm 2 to 0.6 A/cm 2 while the temperature of said solution is being kept at 25°-60° C. In the method of the present invention, the concentration of the hydrogen chloride in the etching solution must be in the range between 0.5 N and 3.5 N. This range is determined by considering it in association with the other factors, in order to form etching pits as many as possible, to distribute these over the aluminum surface as evenly as possible in conformity with the quality of aluminum material dominated particularly, by its purity and the orientation of the crystal grain, and in order to reduce the adverse effect of the oxide film covering the surface of the aluminum foil. When the concentration of hydrogen chloride in the solution is under 0.5 N adverse effect due to the oxide film covering the surface of the aluminum foil and/or the detective parts on the surface of the aluminum foil will become great to result in an uneven distribution and a decrease in the number of the etching pits. On the other hand, a concentration of the hydrogen chloride over 3.5 N may result in the increase in the number of the pits but may create drawbacks of surface dissolution and of too shallow undergrown pits. The additive, incorporated into the etching solution together with the hydrogen chloride as an auxiliary agent capable of anodizing aluminium to form a porous oxide film thereon, has an advantage of inhibiting the dissolution of aluminum surface by virtue of its oxidizing activity on the aluminum metal. Namely, in a spot of the surface of the aluminum foil where an oxide film is formed, the etching action is weakened while another spot of the surface being chemically active is only subjected to the etching action to promote a deep growth of an etching pit. At the same time, by preventing chemical dissolution of the metal surface such as chemical or electrolytic polishing which may occasionally occur at the surface, the effective surface area of the metal is greatly extended. These additives may be exemplified as; aliphatic monocarboxylic acids such as acetic acid and propionic acid; aliphatic dicarboxylic acids such as oxalic acid, malonic acid, succinic acid, glutaric acid, adipic acid, maleic acid and fumaric acid; aliphatic oxcarboxylic acids such as glycolic acid, lactic acid, malic acid, tartaric acid and citric acid; and inorganic acids such as nitric acid, phosphoric acid, sulfuric acid, boric acid and .[.sulfamilic.]. .Iadd.sulfamidic .Iaddend.acid. These acids may be used as such, i.e. free acids or as any salts thereof, and may be used singly or as one combined with one or more species selected from those exemplified above. The total concentration of these additives must be kept in a range between 0.001 mole/liter and 0.4 mole/liter. At a concentration under 0.001 mole/liter, no advantage will result from the addition because it only effective to the increase in the number of the etching pits but no effect is expected in the deep growthes of these pits. At a concentration over 0.4 mole/liter, the effective surface area decreases because the addition is only effective to the deep growthes of the etching pits of the larger diameters but of the smaller numbers. Incidentially, there has hitherto been known the use of etching power supply of direct current in a waveform containing a large amount of ripple component or having been superimposed with alternating current. Any of these power supplies are fundamentally classified as direct current and have a drawback that the etching pits may deeply grow into the metal to penetrate there through to give the foil of larger effective surface area but simultaneously, to extremely lower its mechanical strength. The present inventors have investigated the configurations of the etching pits by observing them with an electron microscope in terms of the influence of the above points. As a result of the observation, they have found that by performing the etching operation under given conditions with an alternating current power supply, a configuration of etching pit, which cannot be expected with a direct current power supply, is obtained and it is possible to grow the etching pits deeply without adversely affecting the mechanical strength of the foil. Namely, when an alternating current is supplied from the power source 1, in an arrangement shown in FIG. 1, through carbon electrodes 2 and 2', the etching of the aluminum foil 5 proceeds as follows. Firstly, when the electric potential of the carbon electrode 2 is positive against that of the carbon electrode 2', each of the chloride ions of negative charge in the etching solution 4 contained in an etching cell 3 moves to the side of the carbon electrode 2, and one of the surface of the aluminum foil 5, facing the carbon electrode 2', is subjected to the etching action of the chloride ions. At the same time, the additive contained in the solution acts to form an oxide film on the surface of the aluminum foil 5. Next, when the potential of the carbon electrode 2 is reversed negative against that of the carbon electrode 2', another surface of the aluminum foil 5 facing the carbon electrode 2 is then subjected to the etching action of the chloride ions because each of the chloride ion moves to the side of the carbon electrode 2'. At that time, although hydrogen gas evolves at the surface of the aluminum foil 5 facing the carbon electrode 2', it serves to disperse the aluminum ions present in the vicinity of the surface of the aluminum foil in a high concentration as well as to make the chemical activity of the surface of the aluminum foil even. Accordingly, when the potential of the carbon electrode 2 is reversed positive again, the etching action will not necessarily be initiated at exactly the same positions of the etching pits formed previously, thus preventing the growthes of straight etching pits which are often observed with the direct current etching process and giving a porous and sponge-like configuration of the surface in a form of an aggregation of minute pits to the foil. Moreover, it is advantageous in that the core portion of the foil remains without being etched to effectively prevent the conspicuous decrease in the mechanical strength, because the etching pit will not grow so deeply to penetrate through the foil as is inevitable in the case of the direct current etching process. Incidentally, in the arrangement shown in FIG. 1, the two carbon electrodes are employed and an alternating current is flown through the electrodes though, it is confirmed that a similar result can be obtained in an experiment wherein an alternating current is flown through one of the carbon electrode and the aluminum foil. In this case, however, the etching action is only effected during period of half cycle of the power supply when a positive potential is applied at the aluminum foil and therefore only half of the supplied electric power is utilized in the etching process. Moreover, the latter is not practical in view point of the equipment, because it necessitates a device for flowing a great amount of current directly through the aluminum foil. In performing the method of the present invention, the frequency of the power source and the current density should carefully be maintained within predetermined optimum ranges. In connection with the frequency of the power source, it is obvious that the lower the frequency, the more the configurations of the etching pits proximate to those obtained with direct current. Particularly with an alternating current of a frequency under 5 Hz, the etching pits grow deeply into the foil but are of large sizes and are not evenly distributed over the surface of the foil. Moreover the number of the pits decrease to reduce the extending rate in the effective surface area. On the other hand, in higher frequency, the conditions of the etching become to demonstrate a fine etched surface peculiar to the alternating current etching. Particularly, when the frequency exceeds 120 Hz, the minute etching pits formed on the surface do not grow deeply into the foil and are limited only to the surface as the surface polishing to reduce the extending rate of the effective surface area as in the case of the frequency being under 5 Hz. In connection with the current density, it is confirmed that when it is low, the etching conditions tend to proximate to those in the directric current etching process to promote a deep growth of the etching pit but to make the distribution of the etching pits over the surface uneven. Particularly with a current density under 0.1 A/cm 2 , the etching pits grow deeply into the foil but are small in number and not evenly distributed over the surface of the foil as in the case of the low frequency power supply. When the current density exceeds 0.6 A/cm 2 , only the surface of the foil dissolve without the deep growth of the pits and no extension of the effective surface area cannot be expected. Time is also a dominating factor which influences the performance of the etching operation as well as the configuration of the etched surface. However this factor is highly dependent on the other factors involved and cannot be determined unconditionally with the above mentioned conditions. What can safely indicated here is that the time should be set to correspond to power consumption approximately equal to .[.170.]. .Iadd.85 .Iaddend.coulomb/cm 2 in the case of foil thickness is 100 μm, or to the dissolution of the aluminum metal into the etching solution up to approximately .[.80.]. .Iadd.8.0 .Iaddend.mg/cm 2 of the foil. In the following description, the present invention will be elucidated in more detail by way of examples, presented in contrast with results of the comparative experiments. EXAMPLES 1 AND 2 (COMPARATIVE EXPERIMENTS 1, 2 AND 3). Under conditions summarized in Table 1 below, the etching treatments on the sample aluminum foils (purity, 99.9%; area, 10 cm 2 , thickness, 100 μm) are performed. TABLE 1__________________________________________________________________________ Aluminum dis- Electrolyte in Frequency of Current Time solved into the etching the power density (min. Temp. the solutionExample solution source (Hz) (A/cm.sup.2) sec.) (°C.) (mg/cm.sup.2)__________________________________________________________________________ HCl 1.3 N1 60 0.5 2.50 45 7.9 H.sub.3 PO.sub.4 0.1 mole/liter HCl 1.3 N2 20 0.3 4.45 35 7.6 H.sub.3 PO.sub.4 0.1 mole/literComparativeExperiment NaCl 3.5 mole/liter1 (D.C.) 0.5 2.30 95 8.0 Na.sub.2 SO.sub.4 0.3 mole/liter HCl 1.6 N HNO.sub.3 0.1 mole/liter2 H.sub.3 PO.sub.4 0.1 mole/liter (D.C.) 0.7 1.45 70 8.1 H.sub.2 SO.sub.4 0.005 mole/liter3 HCl 1.3 N 60 0.5 2.50 45 8.0__________________________________________________________________________ In evaluating the performance of the etching processes, each of the etched sample foils is first anodized in an aqueous electrolytic solution containing boric acid at 3% and ammonium borate at 3%, at 30° C.±2° C., up to the formation of oxide film which corresponds to the final formation current under 0.5 mA/10 cm 2 at 80 V, and its electrostatic capacitance is measured in an aqueous ammonium borate (8%) solution at 30° C. The durability of the etched foil against folding distorsion (hereinafter, to be simply referred to as "bending strength") is measured by repeatedly bending (90°) a free end of the foil (10 mm in width) dangled with a weight of 200 g from a clip whose grasping tips are rounded (radii of curvature, 0.2 mm) until being finally cut to be separated. The count in the repeat of each of the right and leftwards bending is taken as a relative scale for expressing the "bending strength". The results of these evaluations are summarized in Table 2 below. TABLE 2______________________________________ Electrostatic Bending capacitance (μF/cm.sup.2) strength (count)______________________________________Example 1 3.7 68 2 5.5 64Comparative 1 5.1 3experiment 2 2.9 5 3 2.0 72______________________________________ As can be apparent from the results shown in Table 2, the aluminum foils etched by the method in accordance with the present invention are of high electrostatic capacitance and of great mechanical strength as compared with those etched by the methods of the comparative experiments. These results also support the importance of a deliberate selection and combination of the conditions, i.e., current density, frequency and temperature which may cause a difference in the electrostatic capacitance of the etched foil. Thus an ingenious contrivance is required in the selection of one these factors in view of the other which may have an interaction or a mutual dependency on the former. EXAMPLES 3-10 (COMPARATIVE EXPERIMENT 4) Next, the influences of the species of the electrolytes in the etching solution on the electrostatic capacitance and mechanical strength of the sample foils are investigated with other operational conditions listed below to obtain the results of the measurements summarized in Table 3. ______________________________________Conditions:______________________________________Sample aluminum foil: purity 99.99% thickness, 65 μmConcentration of HCl: 2.0 NCurrent density: 0.4 A/cm.sup.2Frequency of power source: 50 HzTime: 2'15"Temperature: 40° C.Aluminum dissolved: 5.0 mg/cm.sup.2into the solution______________________________________ TABLE 3______________________________________ Additive and its Electrostatic concentration in the capacitance at Bending etching solution 80 V formation strengthExample (mole/liter) (μF/cm.sup.2) (count)______________________________________3 Acetic acid (0.35) 3.5 324 Oxalic acid (0.35) 4.1 305 Glycolic acid (0.35) 3.9 346 Sulfuric acid (0.01) 4.4 317 Boric acid (0.15) 3.2 308 Sulfamidic (0.15) 4.0 319 Nitric acid (0.015) 2.9 35 Nitric acid (0.05)10 Phosphoric acid (0.05) 4.7 29 Sulfuric acid (0.005)Comparative experiment4 None 2.3 34______________________________________ As indicated in the above Table 3, an addition of one or more auxiliary agent capable of anodizing aluminum, i.e., forming an oxide film on the surface of the metal, is found to prove a great advantage to the enlargement in the electrostatic capacitance of the sample foil. EXAMPLES 11-14. Next, the influence of the current density at the etching process on the electrostatic capacitance and the mechanical strength of the foil is investigated to obtain results summarized in Table 4 below. The other conditions employed in this investigation are set as follows: Sample aluminum foil: purity, 99% thickness, 90 μm Electrolytes in the aqueous etching solution: HCl: 1.0 N Oxalic acid: 0.06 mole/liter Frequency of the power source: 40 Hz Time: 1'15" Temperature: 40° C. Aluminum dissolved into the solution: 7.0 mg/cm 2 TABLE 4______________________________________ Electrostatic Current density capacitance BendingExample A/cm.sup.2 (at 80 V formation) strength (count)______________________________________11 1.0 3.9 5812 0.3 5.1 5513 0.1 3.7 5314 0.03 2.2 47______________________________________ As indicated in the above Table 4, the electrostatic capacitance of the sample foil is greatly influenced by the current density. Under similar conditions as described above, experiments are performed to obtain results represented by FIGS. 2, 3, 4 and 5, respectively. FIG. 2 illustrates the interaction between the current density and the temperature of the etching solution, which influences on the electrostatic capacitance of the etched foil. FIG. 3 illustrates the changes in the electrostatic capacitances of the etched foils, with the changes in the additives, i.e., sulfuric acid, oxalic acid and glycolic acid, and in their concentration in the etching solution. FIG. 4 illustrates the changes in the electrostatic capacitance of the etched foils, with the changes in the concentration of the hydrogen chloride in the etching solution, and in the temperature of the solution, and FIG. 5 illustrates the changes in the electrostatic capacitance of the etched foils with the changes in the formation voltage in the cases of the frequencies of the etching power source are 20 Hz, 60 Hz and 100 l Hz, respectively.

Optimum conditions, employed in an etching process of an aluminum foil for use in an electrolytic capacitor, are selected and combined together for realizing an extension of the effective surface area of the foil. This extension also gives an enlargement of the electrostatic capacitance of the foil electrode as well as an enhancement of the mechanical strength of the foil itself, the latter had been considered to be hardly compatible with the former.

REFERENCE TO RELATED PATENT APPLICATIONS The present patent application is a division of U.S. application Ser. No. 09/547,395, filed Apr. 11, 2000, now U.S. Pat. No. 7,052,812 which is related to, and claims priority from, U.S. provisional patent application Ser. No. 60/128,989 filed Apr. 12, 1999, for THREE-DIMENSIONAL OPTICAL MEMORY IN FLUORESCENT DYE-DOPED PHOTOPOLYMER to the selfsame inventors as is the present application. The content of the Ser. No. 09/547,395 patent application and the 60/129,989 application is incorporated by reference herein. The present application is related to U.S. patent application Ser. No. 09/547,396, now U.S. Pat. No. 6,501,571 filed on Apr. 11, 2000 for THREE-DIMENSIONAL HOLOGRAPHIC STAMPING OF MULTI-LAYER BIT-ORIENTED NON-LINEAR OPTICAL MEDIA, also to the same inventors. The content of the related patent application is incorporated herein by reference. BACKGROUND OF THE INVENTION 1. Field of the Invention The present invention generally concerns optical data storage, optical memories, polymers, multi-layer optical storage, and volume, or 3-D, optical data storage. The present invention particularly concerns the radiation writing of monolithic multi-layer optical storage media, radiation-written monolithic multi-layer optical storage media, and processes of radiation-induced diffusional re-distribution of molecules within a medium for the purposes of data storage. 2. Description of the Prior Art Because of the relative economy and high storage capacity of optical media, optical data storage in optical media has proven to be a convenient means for the recording and physical distribution of large quantities of information. However, despite the current search for shorter wavelength lasers and more efficient coding techniques, standard planar optical storage technologies may not be able to keep up with the ever-increasing demand for faster data rate, higher capacity, media. As a result various new alternatives to optical technologies are currently being investigated. These alternatives include multi-layer reflective media, near-field recording, volume holography, two-photon storage, bit-oriented index modulation, and others. With the exception of near field recording all of these technologies take advantage of the increase in capacity that can be gained by expanding data storage into the third dimension of a thick media. SUMMARY OF THE INVENTION The present invention contemplates a new monolithic multi-layer optical storage media in which data may be stored through the diffusional re-distribution of fluorescent molecules within a polymer host. In accordance with the present invention, the active portion of a radiation-writable media consists of (i) a photopolymer doped with (ii) a dye, preferably a fluorescent dye. The (i) photopolymer becomes selectively polymerized, and/or the (ii) dye—which can be (but need not be), and preferably is, a photoinitiator of photopolymerization—promotes photopolymerization, where and when radiation, normally light radiation is high, most typically at the focal point of a lens, preferably a high numerical aperture lens. The dye molecules bond to the polymer matrix. As they do they become more highly concentrated in the polymerized regions, resulting in a selectively modulated data pattern, or the storage of information. When the dye is, as is preferable, fluorescent, then the optical readout of stored data is based on detection of fluorescence rather than index modulation as in other photopolymer-based memories. Accordingly the problems of media shrinkage and optical scatter are less of a concern. An intensity threshold observed in the recording response of this media due to the presence of inhibitor molecules in the photopolymer permits the three-dimensional confinement of recorded bits, and therefore multi-layer recording. The non-linear recording characteristics of this material have been (i) investigated through a simple model of photopolymerization and diffusion, and (ii) verified experimentally. Both single and multiple layer recordings have been demonstrated. In comparison to the prior art, where changes in volume media have previously been photochemically induced, the present invention offers several advantages. First, molecules within the pre-polymerized, un-exposed, optical medium exhibit such high diffusion and mobility as is, at the initiation of the process, desirable. This high mobility—as opposed to, for example, a solid—makes that constituent components of the volume media can “move around”. In the prior art this has been useful so that the components can undergo photo-induced chemical reaction, or photochemistry. In the present invention the mobility is more concerned with the movement of a photo-detectable material, i.e., a dye. However, when the prior art showed mobility of the constituent components of a volume media, such as in a liquid, this mobility proved troublesome. The same mobility that promoted significant photochemistry made that chemicals and photochemicals, both unchanged and changed, tended to “wander” in the matrix, “greying out” an exposed optical memory and “smearing” such domains (voxels) as are recorded. However, in the present invention this initial mobility will cease. Initially the dye molecules are drawn to regions of exposure, depleting surrounding regions and increasing contrast. However, over a short period the radiation (light) exposure will cause such spreading polymerization as ultimately locks all constituent components into the positions by then assumed. This is an essential characteristic, and forte, of the media of the present invention: mobility followed by immobility. After such minor lapse of time as permits all chemicals and photochemicals to “move into place” (over molecular dimensions), everything “locks up” and nothing is free to move, essentially forever. This is very beneficial. Effectively nothing changes in a memory once written. The dye can “bleach out” slightly after prolonged exposure to bright light, and/or over millions of radiation interrogations, but is essentially stable, and maintaining of a bright fluorescing capability indefinitely long. The exposed and polymerized radiation memory is essentially just as stable as present (circa 2000) memories of both optical and magnetic types that it serves to replace. 1. A Radiation Recordable Medium In one of its aspects the present invention concerns an initially homogeneous medium that is suitably selectively differentiated by selective exposure to radiation. What the medium consists of is rather straightforward: (i) a binder containing a (ii) photopolymer and (iii) a dye. What transpires within the medium in the event of selective illumination is rather more complex, and requires an explanation based on theory. Although it is well established that theory of operation of an invention is not required for an enabling disclosure, we as inventors hereafter hypothesize and postulate and explain, as best we can, what we think is transpiring. Notably, it is possible in our theories of our invention to place greater emphasis on certain characteristics, and phenomena, in some versions of our theories than in others. For example, and as will become clearer as our theories are expounded, in some versions of our theory the propensity of the preferred dye to be a photoinitiator of photopolymerization can be integrated into the theory of the invention, while in other, preferred, versions of the theory of the invention this property of the dye can be considered merely incidental, and adjunct, to the main process. Accordingly, our theories of our invention are not only hypothetical, but capable of being expressed in different words, and with different emphasis on the various component materials and processes. A practitioner of the photochemical arts reading our following theoretical descriptions, and considering the present invention, will thus realize that there are many different ways to describe the underlying phenomena, and that the scope of the present invention should not be unduly limited by arguments concerning the particular words that are variously chosen to express the various theories of the invention. Instead, the process of selective differentiation occurring in the present invention should be understood to broadly encompass the underlying phenomena, howsoever expressed. 2. A Composition of Matter In a one of its embodiments, the present invention may be shortly, and broadly, be conceived of as a composition of matter selectively permanently differentiable in regions by selective radiation illumination. The composition consists essentially of (i) a monomer photopolymerizable to a polymer (i.e., a photopolymer), and (ii) a dye attaching the polymer. The dye both (i) migrates to, and (ii) is fixed in position at, regions photopolymerized by selective illumination with radiation. Preferably (but not necessarily), the dye chemically attaches the polymer at these regions. The composition may further include a binder in which is present the (i) photopolymer, and (ii) the dye. The dye is preferably fluorescent. 3. A Radiation-Alterable Medium The expression of the invention within the previous section 1. is—while distinguished from the prior art by the (i) migration, and the (ii) fixation, of a dye—is likely so short as to make it difficult to fully understand, and appreciate, the theorized operation of the invention. In somewhat greater detail, the present invention may alternatively be considered to be embodied in an initially homogeneous medium that is suitably selectively differentiated by selective exposure to radiation. The medium consists essentially of (i) a binder, containing (ii) a photopolymer, the photopolymer initially substantially homogeneously doped with (iii) a dye that is photoexcitable to bind to photopolymer. Upon selective exposure of certain regions of the medium by radiation, each of (i) polymerization, (ii) dye migration and (iii) dye fixing will occur. The photopolymerization selectively progressively solidifies the photopolymer. Meanwhile, simultaneously, the dye molecules initially migrate to photopolymerized regions, and are there fixed in position by the photopolymer. By this coaction a concentration gradient of dye in the matrix from unexposed to exposed regions results. In other words, a concentration gradient results from diffusion of the dye from unexposed to exposed areas. The initially homogeneous matrix thus becomes selectively differentiated in that dye concentration is increased in exposed areas relative to unexposed areas. The dye is preferably fluorescent. The dye may be, but need not be, a photoinitiator of the photopolymerization. The dye may be, but need not be, fixed in position by the photopolymer by action of chemically attaching the polymer. (It will in any case be mechanically fixed within the polymer.) In some instances—as is the case with the preferred chemical constituents of the invention—the dye chemically attaching the polymer also chemically attaches the monomer. This makes that the dye must, and does, migrate in accompaniment to its attached monomer until, and where, the monomer does photopolymerize to a polymer. 4. A Radiation Recordable Medium In still yet another of its aspects, the present invention may be considered to be embodied in a medium suitably selectively recorded by radiation. The medium again consists essentially of (i) a host binder, containing (ii) a liquid monomer, in which molecular mobilities are relatively higher, that is photopolymerized into a solid polymer in which molecular mobilities are relatively lower, the monomer initially substantially homogeneously doped with (iii) a dye that is photoexcitable to bind to at least the polymer. In terms of the most detailed theory of which we the inventors are possessed, we believe, and have experimental data supporting, that upon selective exposure of certain areas of the medium by radiation, both (i) polymerization and (ii) dye fixing will occur. The monomer undergoes a polymerization process that solidifies the medium. Meanwhile, simultaneously, radiation-exposed and photoexcited dye molecules bind to at least polymer molecules (and normally to the monomer molecules also, although this is not essential). These dye molecules photoexcited to bind to at least the polymer molecules become, due to the relatively lower molecular mobility in the polymer, relatively fixed in their location, Meanwhile other, un-photoexcited, dye molecules remain relatively more mobile. This results in an initial migration and redistribution of dye to photoexposed regions until, dye migration being substantially complete, photopolymerization occurs, locking the migrated dye in place at a relatively higher concentration at the selectively exposed regions. By this coaction of (i) polymerization and (ii) dye fixing there arises a concentration gradient of dye molecules in the material from the unexposed to the exposed areas. This concentration gradient results, of course, from diffusion of the dye from unexposed to exposed areas. This dye concentration is, of course, increased in exposed areas relative to unexposed areas. After the selective radiation exposure stops an excess concentration of dye molecules in the exposed areas serves as a permanent record of the selective radiation exposure. Regardless of exactly how the composition of matter, or medium, or process of the present invention is expressed, many specific aspects of construction and operation are preferred in implementation of the present invention. Using the particular language of this section 5., the “medium” preferably further includes an inhibitor of the photopolymerization so that in regions of low radiation exposure polymerization is inhibited nonetheless that other regions where radiation is concentrated become fully polymerized. This inhibitor of the photopolymerization typically consists essentially of oxygen. In the most preferred embodiment the binder consists essentially of (i) a matrix, for example cellulose acetate propionate, and (ii) a solvent, for example acetone. Meanwhile, the photopolymer consists essentially of (i) a monomer, for example dipentaerythritol pentaacrylate; (ii) a crosslinker, for example 1-vinyl-2-pyrrolidinone; (iii) an initiator, for example N-phenyl glycine; and (iv) a photosensitizer, for example camphor quinone. A polymer chemist will recognize the function of these constituent components, and the utility of each within a polymerizable system. Meanwhile, the dye is preferably fluorescent, and, although permissively drawn from a large group, preferably consisting essentially of (i) Rhodamine B; and/or (ii) Bodipy Red. In use of the medium the photopolymer is initially substantially uniformly doped with the dye which, as stated, is preferably fluorescent. The medium is selectively illuminated in regions to write and to store data as an optical memory, reading of the optical memory transpiring by inducing fluorescence of the dye. When the selectively illuminated regions are voxels, the optical memory becomes a three-dimensional volume optical memory. Reading of the optical memory transpires by detecting relative presence or absence of the dye, preferably through its fluorescent light emissions when stimulated. 5. A Method of Radiatively Recording Information in a Photosensitive Medium In still yet another of its aspects the present invention is embodied in a method of radiatively recording information in a photosensitive medium. The method commences by establishing a host binder, or matrix, containing (i) a photopolymer that is polymerized by radiation initially doped with (ii) a dye that photoinitiates photopolymerization. Here again is that theory of the invention, one of many, where photoinitiation is promoted by the dye. There is as much support for this theory as any other: the most preferred dye is in fact a photoinitiator of photopolymerization. It is not, however, believed critical to the present invention that it should be a photoinitiator. It is thus, as previously stated, very much a matter of what is emphasized in expounding different, but equally plausible, theories of the present invention. This matrix is then selectively exposed in selected regions by radiation so that the dye migrates to redistribute itself to exposed regions until, dye migration being substantially complete, photopolymerization occurs, therein serving to lock the migrated dye in place at a relatively higher concentration at the selectively exposed regions. The regions where is present the dye represent the information recorded. The recording of information is most commonly for the purpose of making a radiation memory. Reading of the radiation memory transpires by radiatively detecting the relative presence, or absence, of the dye in the regions of the matrix. This radiative detecting of the relative presence, or absence, of the dye is preferably by inducing fluorescence in fluorescent dye. 6. A Radiation Memory System In still yet another of its aspects the present invention is embodied in a radiation memory system that is at least writable, and that is preferably read-writable. The system includes (i) a first laser; illuminating with a first laser beam (ii) a mask; which is imaged through (iii) a 4-f lens system; to a certain depth within (iv) a volume optical recording medium. This medium consists essentially of (i) a host matrix, containing (ii) a photopolymer that is polymerized by radiation, doped with (iii) a dye having a greater affinity for photopolymer than for the host matrix. Upon selective exposure by the first laser through the mask of certain regions of the matrix, the dye will initially migrate to redistribute itself to radiatively exposed regions until, dye migration being substantially complete, photopolymerization will occur, locking the migrated dye in place at a relatively higher concentration at the selectively exposed regions. This writable radiation memory system may optionally, between the first laser and the mask, further include a rotating diffuser serving to break up the spatial coherence of the first laser beam and randomizing the diffraction image of the first laser beam along its propagation path except when it is in focus. It may still further optionally include a lens for bringing to a focus the first laser beam within a volume optical recording medium. The memory system may be expanded and extended to reading. Such a read-write radiation memory system further preferably includes (i) a shutter means for blocking the first laser beam; (ii) a second laser for producing a second laser beam; (iii) a lens for focusing the second laser beam onto selectively exposed regions of the matrix that are at a particular depth within the volume optical recording medium, causing the exposed regions to fluoresce; and (iv) a detector of the fluorescence as representing the read contents of the volume optical recording medium, ergo a read-write optical memory. The (iii) lens particularly preferably consists of a cylindrical lens, and the (iv) detector is preferably a CCD. These and other aspects and attributes of the present invention will become increasingly clear upon reference to the following drawings and accompanying specification. BRIEF DESCRIPTION OF THE DRAWINGS Referring particularly to the drawings for the purpose of illustration only and not to limit the scope of the invention in any way, these illustrations follow: FIG. 1 , consisting of FIGS. 1 a through 1 c , is a diagram showing the mechanism of recording where in (a) the media begins with a uniform distribution of fluorescent dye throughout the volume; in (b) light begins photopolymerization and simultaneously fixes dye to the polymer matrix while free dye will diffuse into the exposed areas leaving the unexposed regions depleted of dye; and in (c) after completion of the polymerization process the dye is held in a permanent polymer matrix. FIG. 2 is a schematic diagram of the optical system used for both recording and readout of the dye-doped photopolymer media. FIG. 3 is a pictorial representation of a resolution target image recorded into the dye-doped photopolymer media and readout by the excitation of fluorescence in the recorded layer. FIG. 4 , consisting of FIGS. 4A through 4D , are pictorial images of the edges of exposed regions and accompanying horizontal profiles where in all cases the left half of the image has first been exposed: FIG. 4A showing edge enhancement observed 30 seconds after exposure as dye begins to diffuse into the exposed region; FIG. 4B five minutes after exposure the profile resembles a flat top while the dip in intensity corresponding to the dye-depleted region is gone; and with FIGS. 4C and 4D showing exposures that have been immediately followed by a uniform fixing exposure so that thirty seconds after the exposure of FIG. 4C the edge enhanced structure is again apparent, however 5 minutes after exposure of FIG. 4D the structure has not changed. FIG. 5 is a graph showing a numerical simulation of component concentrations and modulation ratio for low intensity exposure, I 0 =0.41 W/cm 2 , where polymerization is inhibited. FIG. 6 is a graph showing numerical simulation of component concentrations and modulation ratio for high intensity exposure, I 0 =0.65 W/cm 2 , where, after a short induction period, polymerization begins. FIG. 7 is a graph showing numerical simulation of modulation ratio as a function of recording intensity for images recorded at three different exposure energies. In all cases a sharp threshold is observed followed by a rapid increased in the degree of polymerization. Below the threshold point polymerization is negligible. FIG. 8 is a graph showing fluorescence contrast ratio of images recorded for varying exposure energies at three different constant intensity levels. FIG. 9 is a graph showing fluorescence contrast ratio of images recorded at varying intensity levels for three different constant exposure energies. FIG. 10 , consisting of FIGS. 10 a through 10 d , are pictorial images readout from four layers recorded in the media 70 μm apart in depth. FIG. 11 is a diagrammatic representation of the manner of the recovery of the recording of the images of FIG. 10 by stepping the focal position of the objective lens starting with the deepest image first. FIG. 12 , consisting of FIGS. 12 a through 12 f , are diagrams showing reactions during photopolymerization (fm. ref. 12 ) FIG. 13 is a diagrammatic representation of grey-level mask array for fabricating conical lenses with uniform light in photopolymer. DESCRIPTION OF THE PREFERRED EMBODIMENT Although specific embodiments of the invention will now be described with reference to the drawings, it should be understood that such embodiments are by way of example only and are merely illustrative of but a small number of the many possible specific embodiments to which the principles of the invention may be applied. Various changes and modifications obvious to one skilled in the art to which the invention pertains are deemed to be within the spirit, scope and contemplation of the invention as further defined in the appended claims. 1. A New Monolithic, Multi-layer Optical Media The new monolithic, multi-layer optical media of the present invention has advantages over other new storage technologies in terms of ease of recording and readout, while sharing the same characteristics of very high density, fast access time, and parallel readout capability. The active portion of the media is a photopolymer—similar to that used for other volume optical media—which has in addition been doped with a fluorescent organic dye. Recording occurs during photopolymerization when the distributions of dye become more highly concentrated in the polymerized regions, therefore producing a modulated fluorescence readout pattern. In this architecture recordings are bit-oriented and are performed in a 2D parallel fashion. The readout which is based on excitation of fluorescence is non-destructive and is detected also in 2D parallel. Since index modulation, or changes in the material density, is not necessary for readout in this media some of the problems commonly associated with photopolymer memories such as shrinkage and optical scatter are less of a concern. This material is non-erasable, therefore given its ease of production and low-cost materials this media is best suited for mass distribution ROM applications. In the following sections we will investigate the properties of this new media. In Section 2 we will describe the material composition and the optical setups used for image recording. The fluorescence modulation process which is based on the diffusion of dye within the material will be described and shown experimentally. In Section 3 we will then examine the non-linear recording characteristics that have been observed for this material. Specifically, a threshold in the recording intensity response will be demonstrated that indicates that this material is suitable for multi-layer data storage. The origin of this threshold will be investigated with a simple model and verified through a series of recordings. In Section 4 we will present the results of a four-layer recording and readout demonstration and discuss the limitations of this media. Finally, in section 5 we give the best and most complete theory of which we are presently possessed as to how, why and in what best mode photopolymerization for purposes of three-dimensional optical data storage in fluorescent dye-doped photopolymer in accordance with the present invention proceeds. 2. Materials and Recording The preferred embodiment of optical recording media in accordance with the present invention consists of an inert polymer host in which is suspended both a fluorescent dye and a complete photopolymer system. The photopolymer used here polymerizes in a standard free-radical reaction and consists of the monomer dipentaerythritol pentaacrylate, the crosslinker 1-vinyl-2-pyrrolidinone, the initiator N-phenyl glycine, and the photosensitizer camphor quinone. The recording wavelength of this material is determined by the photosensitizer which in this case has an absorption peak very close to the blue 488 nm line of an argon ion laser. In these results we have chosen Rhodamine B as the fluorescent dye which absorbs well at the 543.5 nm wavelength of a green helium neon laser and has an orange fluorescence band from 560 nm to 640 nm. We have also tried a number of other dyes with similar success including Bodipy Red which has an absorption peak at 650 nm and fluorescence at 665 nm, thus allowing compact red diode lasers to be used for readout. All of these active components are suspended in a binder of cellulose acetate propionate and dissolved in acetone. The pre-polymer mixture must be kept in a semi-viscous state in order for the active components to have a sufficiently high diffusion rate, therefore samples are prepared in sealed square cuvets. After polymerization the material becomes solid and inert. Difficulties may arise from having a pre-polymerized material that is not solid since layers within the volume are potentially free to drift within the viscous medium unless the recorded layers are anchored to the bottom or all of the layers are recorded very rapidly in a simultaneous manner. For these experiments we have used the former approach to avoid movement of the data, however for this material to be useful as a mass distribution ROM media an appropriate optical data stamping technique should be developed so that information can be duplicated in the media in a rapid, inexpensive fashion. Implementing massively parallel recording may be one approach to rapid optical data stamping. The mechanism of recording in this material is shown in FIG. 1 . Before recording the fluorescent dye molecules are evenly distributed throughout the media, however during the recording process a modulation of the fluorescence intensity according to the recorded pattern is developed. We believe that this modulation results from a binding of dye molecules to the polymer matrix since it has been observed that Rhodamine B can act itself as a photoinitiator for this monomer without any other additional dopants. Two processes are therefore initiated by exposure to blue light: polymerization and dye fixing. In the exposed areas, the monomer undergoes a polymerization process that solidifies the matrix, while, simultaneously, photoexcited dye molecules bind to both monomer and polymer molecules. Dye molecules that are incorporated in the polymer matrix will become fixed in their location resulting in a concentration gradient of free dye molecules in the material. This concentration gradient causes more dye to diffuse in from the dark regions thereby increasing the total fluorescence level observed in the exposed areas and reducing that in the unexposed areas. Once the material has been fully polymerized and all of the free dye molecules have been bonded the fluorescence modulation is fixed and recorded images are permanent. Materials that operate on a principle similar to that described here have also recently been developed for holographic memories where the goal there is to maximize index modulation rather than fluorescence modulation. FIG. 2 shows the optical setup used for both recording and readout of data. An argon laser illuminates a mask which is imaged through a 4-f system to a certain depth within the recording media. The rotating diffuser immediately before the mask is necessary to reduce inter-layer crosstalk by breaking up the spatial coherence of the source and randomizing the diffraction image of the beam along its propagation path except when it is in focus. An Olympus 40x/0.6 N.A. microscope objective that can variably compensate the spherical aberrations of up to 2.5 mm of plastic or glass is used to focus into the media. The position of the microscope objective is stepped axially to image to different depths within the media for multi-layer recordings. Images are recorded in exposures typically ranging from 2 sec to 20 sec with the argon laser focused down to recording intensities between 1 and 10 W/cm 2 . Readout is performed by blocking the recording beam and illuminating the recorded layer orthogonally with 350 μW from the green He—Ne. The green beam is cylindrically focused into a thin sheet of light approximately 15 μm thick to select just one layer out of the volume. The fluorescent dye absorbs the green light and then emits an orange fluorescence that is collected and imaged onto the CCD detector. The orthogonal nature of the readout scheme permits 2D parallel detection of data which in a high speed readout system may allow extremely high data transfer rates. A typical single plane recording of an air force resolution target is shown in FIG. 3 . The average contrast ratio of this image is less than 2:1. The smallest resolvable line pair group is group 5:6 which after de-magnification corresponds to 1.28 μm/line pair, or a 0.64 μm minimum feature size. The material therefore exhibits sub-micron resolution which is a necessary requirement for a high density optical storage media. In this picture the green illumination beam is entering from the right side. On the left side of the image it is possible to see some non-uniformities in the readout beam which look like shadows. These shadows are a result of the small index changes that occur when the material polymerizes causing scatter and focusing in the beam. No attempt has been made in this work to minimize the index of refraction modulations, however further optimization of the material composition would be desirable in order to reduce these effects and improve the optical quality of the media. Recorded samples that have been stored in the dark for several months have also shown a slight degradation in image quality most likely from evaporation of the solvent due to leaks in the seal of the cuvet. A more stable solvent would be preferred to improve the shelf life of the material. The diffusion process of the dye during recording can be observed by exposing a large region in the material. FIGS. 4A-4D show four images in which the left half of the sample has been exposed to light. Immediately after exposure an edge enhancement of the boundary between the exposed and unexposed regions may be observed as dye diffuses into the exposed region (see FIG. 4A ). Dye concentration builds up most quickly closest to the edges leaving a dye-depleted region immediately outside the exposed area. Over a period of minutes the entire exposed region will slowly rise to a uniform fluorescence level (see FIG. 4B ). The edge enhanced recording observed in 4 A can be fixed, however, by immediately applying a second exposure that uniformly illuminates the entire image and bonds all of the dye in the sample. FIGS. 4C and 4D show a fixed image immediately after the second exposure and also 5 minutes later. It can be seen that the characteristic edge enhanced structure is still evident even several minutes after the exposure. Because the fluorescent signal strength of a recorded data image is dependent on diffusion of dye from neighboring unexposed regions, it can be seen that the recorded fluorescence intensity may vary depending on the size of the recorded bit and the number of neighboring bits which have been recorded. For recording of digital data modulation codes such as a bi-phase (or Manchester) code may be necessary to ensure uniform signal strength of recorded bits. 3. Non-linear Polymerization 3.1 Modeling In order for any particular media to be viable as a three-dimensional optical media it must be possible to selectively record layers within the volume without adversely affecting the information stored at other layers. This material exhibits an intensity threshold in its recording curve which results from the presence of the inhibitor oxygen. Oxygen is dissolved in the mixture because the samples have been prepared in air. In most cases, the presence of oxygen in photopolymers exhibits itself as an exposure time threshold, however in this material an intensity threshold may also be observed due to the material's high diffusion rate. The existence of this intensity threshold is independent of the presence of the fluorescent dye and can be demonstrated through a simple model of the polymerization process of a photopolymer. The processes that will be considered in this model are: free radical generation by light, chain initiation (the creation of a radical site on a monomer or polymer), chain propagation (the addition of monomers to a polymer chain), chain termination (the neutralization of a polymer radical site), radical annihilation (the neutralization of free radicals with each other), inhibition (the neutralization of radicals by oxygen), and diffusion. We wish to write down the rate equations for the following components of the photopolymer: the monomer (M), the polymer (P), the inhibitor (H), the free-radical initiators (R.), the radical sites on the polymer chain (P.), and the neutral components such as the solvent or binder (N). For some arbitrary input light intensity I(x,y,z,t) the rate equations are given as: ∂ [ R · ] ∂ t = ( ση h ⁢ ⁢ ω ) ⁢ I ⁡ ( x , y , z , t ) ⁡ [ R ] + α R ⁢ ∇ 2 ⁢ [ R · ] - k i ⁡ [ R · ] ⁢ ( [ M ] + [ P ] ) - k a ⁡ [ R · ] 2 - k t ⁡ [ R · ] ⁡ [ P · ] - k h ⁡ [ R · ] ⁡ [ H ] ( 1 ) ∂ [ H ] ∂ t = α H ⁢ ∇ 2 ⁢ [ H ] - k h ⁡ [ H ] ⁢ ( [ R · ] + [ P · ] ) ( 2 ) ∂ [ M ] ∂ t = α M ⁢ ∇ · ( [ N ] ⁢ ∇ [ M ] - [ M ] ⁢ ∇ [ N ] ) - k p ⁡ [ M ] ⁢ ( [ R · ] + [ P · ] ) ( 3 ) ∂ [ P · ] ∂ t = k i ⁡ [ R · ] ⁢ ( [ M ] + [ P ] ) - k t ⁡ [ R · ] ⁡ [ P · ] - k h ⁡ [ P · ] ⁡ [ H ] ( 4 ) ∂ [ P ] ∂ t = k p ⁡ [ M ] ⁢ ( [ R · ] + [ P · ] ) ( 5 ) where the k's correspond to the rate coefficients for the polymerization processes of chain initiation (k i ), chain propagation (k p ), chain termination (k t ), radical annihilation (k a ), and inhibition (k h ), the α's correspond to the diffusion constants for free-radicals (α R ), inhibitors (α H ), and monomers and neutral components (α M ), σ is the absorption cross-section at the given wavelength, and η is the conversion efficiency of the initiator into free-radicals. We assume that the number of free-radicals generated is small compared to the total concentration of initiators (R). The various constants in these equations have either been measured from the material or taken as typical values from the literature. Polymer and radical sites on the polymer are assumed not to diffuse. The remaining components must inter-diffuse as a multi-component system. In order to simplify this problem we may observe that the concentration of free-radicals and inhibitors is much smaller than that of monomers, polymers, and neutral components and therefore may be approximated to diffuse freely. We are left then with diffusion in a bi-component system (between monomers and neutral components) which—if we assume the densities of these components are similar—may be approximated as shown in Eqn. 3 subject to the conservation of mass equation: [ M]+[P]+[N ]=constant  (6) As polymerization occurs the concentration of monomer will drop, leading to diffusion of monomers into the exposed areas and expulsion of the neutral component. Polymerization will cease when the neutral component has been completely expelled from the irradiated region and all of the monomer has been converted to polymer. As a measure of the degree of both the polymerization and diffusion in the material we will define a modulation ratio which is directly proportional to the polymer plus the monomer content: modulation ⁢ ⁢ ratio = [ M ] ⁢ ( 0 , t ) + [ P ] ⁢ ( 0 , t ) [ M ] ⁢ ( 2 ⁢ a , t ) + [ P ] ⁢ ( 2 ⁢ a , t ) ( 7 ) where the modulation is compared between the center of the exposed region and a point one diameter away. A high modulation ratio will indicate a high degree of polymerization and diffusion of the monomer. We wish to simulate the recording of a single bit in the media. To reduce the complexity of the computation a spherically symmetric intensity distribution was chosen, I ⁡ ( r ) = { I 0 , for ⁢ ⁢ r ≤ a 0 , for ⁢ ⁢ r > a ( 8 ) allowing the solution to be solved just in terms of r and t. The rate equations are solved numerically for a given set of parameters. FIGS. 5 and 6 show the component concentrations at r=0 during an exposure for two different intensity levels. We can see that for a low recording intensity, I 0 =0.41 W/cm 2 (see FIG. 5 ), the rate of radical generation and inhibitor consumption is slower than the diffusional rate of the inhibitor and therefore the inhibitor concentration never reaches a zero level but instead plateaus into a quasi-equilibrium state. Polymerization in the exposed region is therefore negligible even for long exposure times. However for a higher recording intensity, I 0 =0.65 W/cm 2 (see FIG. 6 ), the inhibitor consumption rate is higher than the diffusional rate and the inhibitor concentration falls to a low enough level for polymerization to occur. It can be seen that after a short induction time monomers begin to be converted into polymer while more monomer flows into the exposed region, resulting in an increasing modulation ratio. Somewhere between the two intensities shown in these two figures there will exist a certain cross-over threshold when the inhibitor consumption rate surpasses the compensating diffusion rate. In FIG. 7 modulation ratio is plotted as a function of recording intensity for three different exposure energies. The exposure time is compensated in each according to the recording intensity to maintain a constant exposure energy. As expected, in each curve there is a distinct intensity threshold followed by a rapid increase in the degree of polymerization. For much higher intensities (and shorter exposure times) the modulation ratio begins to fall off again since there is less time for monomer components to diffuse. Radical annihilation—which is proportional to the square of the radical concentration if we assume that annihilation occurs as the re-combination of two free-radicals—will also become more of a factor at higher intensities and also reduce the polymerization rate. In designing a material for a three-dimensional media, the sharp threshold observed here is the desired characteristic. For a beam sharply focused into the material, recording will occur only at the most highly focused spot where the intensity reaches above the recording threshold. Elsewhere, the media will not be affected. Sharpening the threshold—making the ratio between the threshold and saturation intensities closer to 1—will have the desirable effect of increasing the number of layers that may be recorded in a single media. The position of the threshold level may be modified by adjusting the balance between inhibitor consumption and diffusion either by increasing the initial concentration of inhibitors through the addition of another inhibitor besides oxygen, or by increasing the diffusion rate in the material through changes in the chemical composition or by heating of the material. The position of the saturation level is determined by the polymerization rate which can also be modified by a number of different parameters. Using the model we are able to choose intelligent directions for optimization of the media that maximize the threshold/saturation ratio. It should be noted that this model excludes many higher-order polymerization effects. For example, in the model we have not taken into account the expected reduction in diffusion rate that will occur as the exposed regions reach higher levels of polymerization. Known as the gel effect, this reduction in diffusion would be expected to both increase the rate of polymerization (as the rate of chain termination is slowed) while decreasing the flow of monomer into the exposed regions. Differences in our model and our experimental results may be accounted for by the gel effect as well as other thermal and optical photopolymerization effects that have not been included. 3.2 Experimental In the model described above we have only examined the polymerization process of a standard photopolymer and we have not included the fixing and diffusion of the fluorescent dye included in the fluorescent photopolymer media. We will now experimentally investigate the fluorescence contrast ratio that results during recording in this media which includes the non-linear processes of both the polymerization and dye diffusion. Although polymerization modulation ratio and fluorescence contrast ratio are not equivalent measurements, we can use the more easily observable fluorescence contrast ratio as some indication of the degree of polymerization since fluorescence modulation will not occur unless the dye position has become fixed. We have experimentally investigated the recording characteristics of this material through a series of recordings under different exposure conditions while measuring the observed fluorescence contrast ratio for 4 μm wide bars. FIG. 8 shows the resulting contrast ratio of recordings for varying exposure energies at three different recording intensities. The two higher intensities show a short induction period followed by polymerization, however for the lowest recording intensity (0.42 W/cm 2 ) no image was formed regardless of the exposure energy. The apparent recording intensity threshold may be observed more clearly in FIG. 9 where contrast ratio is plotted over a range of different intensities. In all recordings the exposure time was adjusted to produce a constant total exposure energy. It can be seen that below a recording intensity of approximately 0.6 W/cm 2 no image is formed. The contrast ratio increases rapidly for intensities higher than the threshold and then plateaus at different levels depending on the total exposure energy. Experimentally we observe the same intensity threshold effect predicted by the polymerization model. This threshold now allows us to perform multi-layer recording. 4. Multi-Layer Recording To demonstrate the multi-layer capability of this media four images at varying depths were recorded (see FIG. 10 ). Images were recorded in order from deepest to shallowest with a spacing between layers of 70 μm. It can be seen that the deeper images are somewhat brighter than the images closer to the surface. We have determined from experiment that this is not a result of the recording order but more likely a difference in the material component concentrations in the thin layer closest (<100 μm) to the surface of the media. Evaporation of the solvent during sample preparation may result in reduced dye mobility and therefore a reduced image contrast. A small degree of inter-layer crosstalk is also visible between the images. Because the illumination beam for the mask was not completely uniform the recording intensity had to be increased in order to record the outer edges of each image. The higher intensity in the center of each image exceeds the intensity threshold condition at neighboring planes resulting in crosstalk. This difficulty highlights the importance of having a very uniform illumination beam maintained at just above the threshold recording intensity. The minimum acceptable spacing between layers is determined by the effective numerical aperture of the recording optics as well as the degree of non-linearity in the recording process. The depth over which an image is recorded will be determined by the rate at which the focused image expands and falls below the recording intensity threshold. The focal depth of a lens is inversely proportional to the square of the N.A. In these recordings the N.A. of the recording optics is only moderately high, moreover the features sizes in the images are very large. When the feature sizes are reduced, diffraction out of the focal plane will more quickly blur the image data and produce a uniform, low intensity background that will not be recorded on preceding or subsequent layers. It should also be noted that when using high N.A. optics it becomes necessary to compensate for the spherical aberrations induced by imaging through a thick medium. This can be done either through the use of aberration corrected optics—as was used here—or fluid immersion lenses which have the added advantage of allowing even higher numerical apertures although typically shorter working distances. In addition to the layer spacing, the total usable thickness of the material will determine the maximum number of layers recordable in a single media. As we have mentioned, the maximum thickness of the media can be limited by optical clarity concerns. Optimizations of the material composition that minimize index of refraction modulations and density changes of the material during recording will allow the use of thicker and therefore higher capacity media. 4. Photopolymerization Theory 4.1 The Photopolymerization Process The process of photopolymerization has been extensively studied, so the reader is referred to other sources for a more exhaustive review. There are in general two types of photopolymerization, free radical and ionic, the difference between the two being that in the latter the photon excites an electron on the monomer directly to a π-electronic state, rather than by a photoinitiator intermediary. This usually requires UV wavelengths, whereupon the molecule de-excites into a radical state, becoming a radical monomer. However, most polymers are highly absorptive in the UV, thus for using materials with a mixture of properties or for creating structures that depend upon the intensity profile, this sort of addition polymerization is not appropriate. Thus embodiments rely upon free radical polymerization to create structures in photopolymer. The basic chemical constituents for radical polymerization reaction are a photoinitiator, a monomer basic unit, and a photosensitizer. The latter determines the wavelength range of photopolymer, since the photosensitizer molecule commences the reaction by absorbing a photon and thus exciting one of its electrons to a higher energy state. If this molecule collides with an initiator molecule within its excitation lifetime, then the initiator molecule transfers an electron to the sensitizer, causing the initiator to become a free radical. When this free radical encounters a monomer, the affinity for the radical is such that a double carbon bond on the monomer is changed into a single bond, allowing the radical to attach itself to the monomer and also leaving a radical site on the monomer unit. Thus there is now left what is called a ‘radical monomer’, which can act as a free radical photoinitiator does, attaching additional monomer units to itself in a serial fashion, while still maintaining a radical site on the chain. This process is called chain propagation, and continues until the chain is terminated by collision with another free radical of the opposite sign. Since this is an random process, different chain lengths are generated, the distribution of the lengths dependant upon temperature, initiator and photosensitizer concentration, etc. and is Maxwellian in nature. One can also add an additional agent called a cross-linker to the above chemical constituents. The cross-link agent creates additional radical sites on the monomer, causing the chain to grow ‘branches’ which increase the molecular weight of the resulting polymer, causing the polymerized region to become denser. Since density changes are correlated with refractive-index changes through the Lorentz-Lorentz equation, the addition of a cross-linker agent is often desirable. Inhibitors such as oxygen are present in any photopolymer mixture that is not prepared in a neutral environment, such as nitrogen atmosphere or vacuum. These inhibitors usually have a higher affinity for the radical initiators than the monomers do, combining with the radical initiators to create a neutral molecule. Thus the photopolymerization process described above does not start until all the inhibitors in a local region that is being irradiated have been consumed. This is manifested on a macroscopic scale as a ‘threshold’ in recording time before reaching the linear recording regime. Although pre-exposure of the photopolymer mixture removes the majority of these inhibitors, continual diffusion from the air into the mixture will always require some intensity threshold to be crossed before polymerization can occur. Consider the localization of the polymerization process. Within a constant irradiance region, there will be concentration fluctuations of initiator, inhibitors, photosensitizers and monomers. These will create localized regions where the recording threshold is crossed first and polymerized ‘seeds’ are formed. Thus in this way concentration gradients (chemical potential gradients, in the language of statistical mechanics) are formed, drawing in more monomer so that polymerization grows outwards from these seeds, while less-active components are forced out. Now consider the polymerization process between adjoining regions, one with a high irradiance and one with a low irradiance. In the former, radical initiators and monomers are being formed at a high rate as compared to the latter. Monomer units are drawn in from the low-illuminated area to the high illuminated area, the monomer traveling on the order of a diffusion length. Of the higher-order effects in the previously described process that couple through the polymerization rate equations to create nonlinear effects between differently illuminated regions, two will be mentioned here. The first comes about because the densification of the polymer as time goes forward causes a larger decrease in the diffusion rate of the photopolymer in regions of high irradiance as compared to those exposed to less intensity. Called the gel effect, this manifests itself as an increase in the polymerization rate since the rate of chain termination is slowed (large molecules become ‘trapped’ and additional radicals find it harder to diffuse past denser chain networks) and the flow of monomer into these regions decreases as well. Combined, one should see a more rapid onset of polymerization than expected while having a lower-than expected index modulation. Additionally, radical-radical collisions leading to radical annihilation are more likely to occur in the high-irradiated regions as compared to the low-illuminated regions. This reaction pathway by removing two radicals decreases the number of radicals that could cause early polymer chain termination, leading to longer, more crosslinked polymer chains. In this way a greater index modulation is set up between the regions of high irradiance as compared to the low irradiance regions despite eventual total polymerization. This is seen in FIG. 4 , showing a one-component formatted polymer block in which a striped irradiance pattern created the high-contrast formatted regions while the low-contrast regions where polymerized later with an uniform irradiance. To further increase this modulation, it is possible to use multiple types of monomers, each with a different density and reactivity ratio, called a c/e ratio. This ratio is a measure of the chemical affinity of one monomer with respect to another monomer for attaching onto a polymer chain. Thus to maximize the index break it is best to select a high molecular-weight (high refractive index) monomer with a high c/e ratio and a low molecular weight monomer with a low c/e ratio. There are two problems associated with polymerization as a method to manufacture precision structures. The first is that the densification of the media causes a volume change which causes shrinkage (change in the physical dimensions of the structure) Thus precise placement of small structures is difficult in polymers which have large shrinkage. However, other inventors and authors have proposed solutions leading to a low-shrinkage polymer. The second problem with using photons as a means for polymerization is light scattering. The seed regions where polymerization commences more often than not grow non-uniformly, creating scattering centers upon termination of polymerization in the structure. Methods to avoid this problem are currently under investigation. 5. Conclusion The present specification has taught a new optical recording media that records information through the re-distribution of a fluorescent dye within a polymer host. This material exhibits an intensity threshold response due to the diffusional nature of the recording process allowing the three-dimensional confinement of recorded bits. Our conceptual description of the recording process has been verified by a simple model of photopolymerization and diffusion, and recordings in this material have been shown to exhibit both high resolution and the capability for multiple layers. Four layers have been recorded on 70 μm layer spacings with only a minimal amount of inter-layer crosstalk being visible. The recording and readout architectures of this memory both permit parallel optical transfer of data which eventually may provide extremely high data transfer rates, while the multi-layer capability of this media indicates that high capacities may be achieved while still allowing fast access times. Improvements must still be made in the material composition particularly in terms of material stability, nonlinearity, and optical quality. Provided that a suitable optical data stamping technique can be developed, the high recording sensitivity and ease of manufacture demonstrated here indicate that this material may have potential as a high capacity, low cost media for information storage and distribution. In accordance with the preceding explanation, variations and adaptations of three-dimensional optical data storage in a fluorescent dye-doped photopolymer in accordance with the present invention will suggest themselves to a practitioner of the optical storage, and optical storage materials, arts. In accordance with these and other possible variations and adaptations of the present invention, the scope of the invention should be determined in accordance with the following claims, only, and not solely in accordance with that embodiment within which the invention has been taught.

A host matrix—normally a binder such as cellulose acetate propionate in a solvent such as acetone—contains a radiation-polymerizable photopolymer—normally a monomer like dipentaerythritol pentaacrylate in combination with a crosslinker like 1-vinyl-2-pyrrolidinone, an initiator like N-phenyl glycine, and/or a photosensitizer like camphor quinine—that is initially uniformly doped with a stable dye—typically Rhodamine B and/or Bodipy Red—that photoinitiates photopolymerization. Upon selective exposure of certain regions of the matrix by radiation, most normally laser light radiation, the dye will initially migrate and redistribute itself to radiatively-exposed regions until, dye migration being substantially complete, photopolymerization will occur, locking the migrated dye in place at a relatively higher concentration at the selectively exposed regions. The dye therein stably located can be optically detected by, preferably, light-radiation-stimulated fluorescence. The medium thus serves as an optical memory, including of the volume type, that can be reliably permanently written quickly and inexpensively at high density.

PRIOR APPLICATION [0001] This application claims priority from and is a divisional application of U.S. patent application Ser. No. 11/034,527 filed Jan. 12, 2005. BACKGROUND OF THE INVENTION [0002] 1. Introduction [0003] This invention relates to an emulsion composition for the formation of an artificial tear film over the ocular surface of the eye capable of providing mechanical lubrication while reducing evaporation of fluid. The composition is also useful for delivering medication to the ocular surface and for treating individuals wearing ocular prostheses such as contact lenses as the composition wets and provides lubrication for both the ocular surface and the surface of the prosthesis. More particularly, the invention relates to emulsion compositions capable of augmenting and maintaining a stable tear film over the ocular surface and/or delivering a medication to said surface without causing substantial blurring of vision nor discomfort. The emulsion is desirably in the form of a meta stable emulsion and is characterized by the use of a surfactant combination suitable for formation such an emulsion and maintaining the integrity of the emulsion during high temperature autoclaving. [0004] 2. Description of the Prior Art [0005] It is known in the art that an aqueous tear film extends over the ocular surface and maintains the ocular surface moist and lubricated. It is also known that dehydration of moisture from the eye may result in discomfort. Further, it is known that compositions are available in the market intended for dry eye treatment. Commercially available compositions are primarily aqueous materials that supplement the tear film by adding a film of a water-soluble polymer over the surface of the eye. This film is short lived and provides limited relief. [0006] The feeling of discomfort resulting from a dry eye condition may include ocular dryness, grittiness, burning, soreness or scratching, dependent upon the subject and the condition of the subject. Proposed causes for dry eye, treatment, and symptoms are described in a compendium of papers edited by Holly, The Preocular Tear Film in Health, Disease, and Contact Lens Wear, The Dry Eye Institute, Lubbock, Tex. 1986; edited by David A. Sullivan, Lacrimal Gland, Tear Film, and Dry Eye Syndromes, 1994, Plenum Press, New York; edited by David A. Sullivan et. al, Lacrimal Gland, Tear Film, and Dry Eye Syndromes 2, 1998, Plenum Press, New York; edited by David A. Sullivan et. al, Lacrimal Gland, Tear Film, and Dry Eye Syndromes 3, Part A and B, 2002, Kluwer Academic/Plenum Publishers, New York incorporated herein by reference for their teachings of the dry eye condition and the treatment thereof. [0007] The most common treatment for dry eye involves temporary alleviation of dry eye symptoms by topical application of a tear substitute that adds a large volume of liquid to the anterior surface of the eye and related adnexa. Typical commercially available tear substitute compositions comprise water-soluble polymer solutions. Examples of such solutions include saline solutions of polyvinyl alcohol, hydroxypropylmethyl cellulose, or carboxymethyl cellulose. U.S. Pat. No. 4,421,748 teaches an artificial tear composition comprising an aqueous hypotonic solution of lecithin and a viscosity-adjusting agent such as a solution of a soluble cellulose. [0008] Methods used to quantify the effectiveness of tear substitutes for dry eye treatment solutions have not been standardized, and many methods used to quantify the results obtained using such tear substitute compositions are often inaccurate. For this reason, it is known that reported relief of dry eye symptoms using known tear substitutes varies considerably from subject to subject, and regardless of the method used to quantify relief using a tear substitute, relief often does not exceed several minutes. [0009] The symptoms associated with dry eye are often exacerbated with subjects using ocular prostheses such as contact lenses. In some cases, contact lens intolerance is caused in part, or in total, by the condition of dry eye and its symptoms. Further, the rate of evaporation from the eye is accelerated by the nature of the contact lens surface and the physical presence of the contact lens results in meniscii formation with additional physical and evaporative effects, even with subjects having an adequate tear film. For many subjects, contact lens intolerance is not overcome by topical application of tear substitutes. Therefore, there is a need for improved compositions and processes for treatment of the dry eye condition and for improving tolerance to ocular prostheses. [0010] Improved compositions for dry eye treatment are disclosed in U.S. Pat. Nos. 4,914,088; 5,278,151; 5,294,607; 5,578,586, each incorporated herein by reference for its teaching of how to form an oil film over the surface of the eye including compositions used therefor. U.S. Pat. No. 4,914,088 teaches the use of certain charged phospholipids for the treatment of dry eye symptoms. The addition of a charged phospholipid to the eye is believed to assist in replicating the tear film that would naturally occur in the eye. In accordance with the patent, the phospholipid composition, preferably in the form of an aqueous emulsion, is topically applied to the eye where it is believed to disperse over the ocular surface and form a film that replicates a lipid layer that would be formed by the spreading of a naturally occurring lipid secreted principally from the meibomian glands during blinking. Because the phospholipid, when applied to the eye, in one embodiment, carries a net negative charge, it is believed that aligned molecules repel each other preventing complex aggregate formation thereby resulting in a stable phospholipid film. The patent theorizes that the film formed from the charged phospholipid assists in the formation of a barrier film reducing evaporation of the aqueous layer, thereby preserving the tear film. It is also now theorized that the phospholipid also functioned as a surfactant maintaining the emulsion stable. [0011] The above referenced U.S. Pat. Nos. 5,278,151; 5,294,607; and 5,578,586 disclose further improvements in dry eye treatment. In accordance with the disclosure of said patents, the dry eye treatment composition of U.S. Pat. No. 4,914,088 is improved by the addition of an oil to the eye treatment composition, preferably a non-polar oil. The oil is added to improve the performance of a dry eye treatment composition by increasing the longevity of the tear film formed on the eye as a consequence of the formation of an oil film over the ocular surface that functions as an evaporation barrier—i.e, by providing and/or thickening the dehydration barrier (the oil layer) on the outer surface of the tear film. Thus, the oil increases the efficacy of the dry eye treatment solution and reduces performance variability from subject to subject. [0012] A preferred embodiment disclosed in the above referenced patents is a dry eye treatment composition comprising a meta stable oil in water emulsion where the water phase includes the charged phospholipid believed to function both as an emulsifier and as a surfactant that assists in spreading of the oil over the eye to form a non-blurring film bonding of the oil to the ocular surface. Preferably, the oil phase comprises a non-polar oil. In accordance with this preferred embodiment, the emulsion is desirably “meta” stable so that when the emulsion is applied to the eye, it will rapidly break and spread over the ocular surface when it first comes into contact with the ocular surface, all as explained in the aforesaid patents. [0013] The meta stable emulsions of the foregoing patents are formulated whereby the total amount of oil added to the eye preferably does not exceed 25 μl, more preferably varies between about 1 and 10 μl and most preferably varies between about 1 and 5 μl. If the amount of oil added to the eye is in excess of 25 μl, the oil layer on the surface of the eye may be of excessive thickness resulting in formation of oil globules on the surface of the eye. These globules are likely to result in prolonged blurring. To achieve control of the amount of oil added to the eye, the concentration limits of the oil in the emulsion are controlled within reasonable limits. An emulsion containing the oil in a concentration of at least 0.1 percent by weight of the total composition provides some benefits, a preferred concentration is at least 1.0 percent of the weight of the treatment composition, and the most preferred oil content varies between about 2.5 and 12.5 percent by weight of the emulsion. [0014] Though the use of an oil in water meta stable emulsion having a negatively charged phospholipid as a surfactant provides excellent clinical results for dry eye treatment, there are certain disadvantages associated with their use. For example, the phospholipid component is costly when manufactured to the requirements and tolerances required for use on the eye. In addition, the storage of the phospholipids requires special conditions. Further, the lack of a long history relating to the use of a phospholipid on the eye could raise questions regarding safety and might create possible concerns by regulatory agencies that might require lengthy and costly clinical trials for approval. A further problem involves possible reluctance of companies marketing eye treatment products to deviate from the use of those ingredients having a long history of uneventful use in existing, commercially available treatment products. [0015] For the foregoing reasons, it is desirable to find one or more surfactants that may be substituted for the charged phospholipids used to form a meta stable oil in water emulsion as disclosed in the aforesaid patents. Though it might appear that simple trial and error could be used to find a suitable surfactant, the task of finding a substitute surfactant is difficult. For example, the replacement surfactant must be acceptable for human use. Many available surfactants are not approved for use on the ocular surface. The replacement surfactant must not cause discomfort to the patient when used in a concentration adequate to form the desired emulsion. Many surfactants may not be added to the eye in suitable concentration without causing stinging. A physiologicall pH of between about 7.0 and 7.8 is required for application to the ocular surface. Many surfactants function as surfactants within a prescribed range of pH, both above and below pH 7. The desired emulsion for treatment of dry eye is preferably meta stable enabling it to rapidly break when applied to the eye. Therefore, the replacement surfactant must enable formation of an emulsion that is stable in manufacture and storage and meta stable and capable of breaking when applied to the ocular surface. The replacement surfactant must be capable of forming an emulsion containing oil in an acceptable concentration as described above to avoid prolonged blurring following application. Finally, the emulsion formed must be sufficiently robust to withstand sterilization at elevated temperatures without breaking, but sufficiently unstable so as to break when applied to the eye. It has been found that many replacement surfactants capable of forming a stable emulsion are incapable of maintaining stability of the emulsion during autoclaving at that temperature required for sterilization if used in a concentration suitable for addition to the eye without causing stinging, or in the alternative, if sufficient to withstand autoclaving, may be so robust that they will not break when applied to the eye. SUMMARY OF THE INVENTION [0016] In accordance with the subject invention, it has been found that a preferred meta stable oil in water emulsion suitable for application to the ocular surface for treatment of the eye may be formed using a combination of surfactants as emulsifiers where one surfactant is a physiologically acceptable surfactant capable of forming the desired meta stable emulsion at physiological pH, hereafter the “primary surfactant”, and an additional surfactant, used in combination with the primary surfactant, is a physiologically acceptable surfactant capable of maintaining the emulsion stable during autoclaving at temperatures in excess of 75° C. or higher without preventing the emulsion from breaking when applied to the eye, hereafter the “secondary surfactant”. [0017] The preferred primary surfactant comprises any one or more physiologically acceptable surfactants capable of forming a meta stable oil in water emulsion at pH between about 7.0 and 7.8 without causing discomfort to the patient when used in a concentration adequate to form the desired emulsion having an oil phase in a concentration of from 1.0 percent by weight up to that amount below that which would causes blurring. The term “meta stable emulsion” means one that is stable in storage but breaks rapidly when instilled onto the ocular surface as described in the above referenced U.S. Pat. Nos. 5,278,151; 5,294,607; 5,578,586. The primary surfactant may be identified by routine experimentation using procedures described below. Surprisingly, other than the phospholipids, the subject of the above referenced patents, no single surfactant has been found capable of use as a sole surfactant to form a meta stable emulsion meeting the guidelines set forth herein though it should be understood that such an emulsion might be formed using a single surfactant in high a concentration whereby the patient is likely to experience stinging when the emulsion is added to the eye. [0018] The preferred secondary surfactant is one or more physiologically acceptable surfactants that is used in conjunction with the primary surfactant which does not alter the meta stable form of the emulsion and does not cause discomfort to the patient at efficacious concentration, while stabilizing the emulsion by preventing decomposition at the elevated temperatures required for autoclaving, typically at temperatures in excess of 75° C. and desirably at temperatures at or in excess of 100° C. Though not mandatory for all secondary surfactants, as a guideline only, the secondary surfactant desirably has a relatively long chain with a minimum of 6 hydrophilic groups and an HLB of 9 or more, and preferably an HLB ranging between 12 and 20, and a lipophilic group that is small in relation to the hydrophilic group and preferably, the same or similar in structure to the lipophilic group of the primary surfactant. [0019] From the literature, it would be expected that one skilled in the art would select a surfactant combination having an arithmetic mean HLB of between about 8 and 14 and more typically, between about 10 and 12 for formation of an oil in water emulsion of the type described herein. The arithmetic mean is determined based upon the HLB of the individual surfactants selected and the concentration of each surfactant used. Unexpectedly, an arithmetic HLB of between 8 and 14 is not required for purposes of the present invention as will be demonstrated below. DESCRIPTION OF THE PREFERRED EMBODIMENTS [0020] The treatment composition of the invention is an oil in water emulsion having an aqueous phase, an oil phase, and a surfactant combination used for the dual purpose of stabilizing the emulsion and spreading the emulsion over the ocular surface following its application to the eye. The surfactant combination comprises a primary surfactant and secondary surfactant and is one that enables formation of an emulsion that is stable in manufacture and during storage, but desirably meta stable when applied to the ocular surface—i.e., one that rapidly differentiates when applied to the eye whereby a non blurring film of oil is rapidly formed over the ocular surface. A stable emulsion during manufacture and storage is one that may separate into separate phases during standing, but can be reconstituted by simple shaking. An unstable emulsion is one that breaks typically forming an oil film or slick that cannot be eliminated by simple shaking. [0021] A meta stable emulsion during use is desirable for purposes of this invention. Though useable for alleviation of dry eye symptoms, a stable emulsion, as opposed to a meta stable emulsion, will not differentiate rapidly when applied to the ocular surface. This is undesirable for the following reasons. An emulsion is typically optically opaque due to the presence of two distinct phases. Therefore, an opaque emulsion over the surface of the eye is likely to cause blurring. The duration of blur is dependent upon the time required for the emulsion to differentiate and form separate layers replicating a tear film. In addition, the emulsion is most easily added to the eye as a standard drop from an eyedropper. The eye is capable of holding a limited volume of fluid—a volume that is less than 25 μl. A volume of 25 μl is substantially less than the volume of a standard drop. Therefore, if the emulsion is stable and fails to differentiate rapidly following application to the eye, excess emulsion will be discharged from the eye during blinking. Discharge of the emulsion from the eye will result in discharge of efficacious components of the treatment solution from the eye before a long lasting tear film can be formed. For this reason, efficacious components may not be available in sufficient quantity to form the desired tear film. Consequently, though a stable emulsion might alleviate the symptoms of dry eye for a limited period of time, it is a lesser preferred embodiment of the invention. [0022] A meta stable emulsion, as the term is used herein, is one that is either stable in storage, or differentiated into two separate layers, but is readily reconstituted by simple shaking prior to use. When a meta stable emulsion is added to the eye as a standard drop, it quickly differentiates permitting rapid formation of an oil film over the corneal surface without excessive oil discharge by blinking. Preferably, the emulsion will differentiate within about 5 blinks following application to the eye, more preferably in a time of less than about 30 seconds. Blurring may occur during the time required to move the bulk of the excess liquid to the canthi and discharge the same from the eye. During and following differentiation of the emulsion, the formation of the oil film is assisted by use of the surfactant combination which serves to help form the emulsion and facilitate the spread the oil over the surface of the eye as the emulsion breaks. Consequently, a meta stable emulsion is the preferred embodiment of this invention. [0023] The surfactant combination used to form a meta stable emulsion must be carefully selected and must meet the following criteria: a. the surfactant combination must enable formation of an emulsion having long term stability, especially when exposed to the high temperatures of autoclaving needed to sterilize the formulation during manufacture, while permitting rapid phase differentiation when applied to the surface of the eye; b. each component of the surfactant combination must be compatible with other components of the emulsion composition and permit formation of the emulsion at the physiological pH of between about 6.5 and 7.8 and preferably, at pH of between 7.2 and 7.5; and c. each component of the surfactant combination must be pharmaceutically acceptable for use on the eye and must be compatible with the eye—i.e., each should be non-toxic and should not cause discomfort such as stinging in the concentrations used. [0027] As described above, the surfactants used to form the emulsions of the invention comprise a combination of a primary surfactant and a secondary surfactant. [0028] The primary surfactant is any one or more pharmaceutically acceptable surfactants that meets the above criteria and desirably forms a meta stable emulsion by itself or in combination with the secondary surfactant, but differs in chemical structure from the secondary surfactant. The literature is replete with thousands of surfactants having a variety of chemical structures described as useful for the formation and stabilization of an oil in water emulsion. To provide an exhaustive list of representative surfactants capable of functioning as a primary surfactant for purposes of the subject invention would be laborious and would omit many useful candidate materials. Therefore, in addition to the representative examples given below, a procedure is given intended to enable one skilled in the art to determine if a given surfactant may be used as a primary surfactant in accordance with the preferred embodiment of the subject invention. This procedure involves the following steps: a. select a surfactant approved for use on the ocular surface within a useful concentration range as given below; b. from the literature or by testing, determine if the surfactant is capable of forming an emulsion with the oil and water components at physiological pH; c. prepare an emulsion having the concentration of emulsion components given below and determine if the emulsion is stable during storage, a minimum of three months under normal storage conditions, or capable of being reconstituted by simple shaking; and d. apply the emulsion to the ocular surface and determine if the emulsion breaks on the ocular surface within a minute or less, preferably in less than 30 seconds or 5 blinks. [0033] Representative examples of primary surfactants meeting the criteria given above include ionic and non-ionic surfactants but non-ionic surfactants are preferred as they are less prone to cause stinging when applied to the eye. Specific examples of the nonionic surfactant include alkyl ethers such as polyoxyethylene octyl ether, polyoxyethylene lauryl ether, polyoxyethylene stearyl ether and polyoxyethylene oleyl ether; alkyl phenyl ethers such as polyoxyethylene octylphenyl ether and polyoxyethylene nonylphenyl ether; alkylesters such as polyoxyethylene laurate, polyoxyethylene stearate and polyoxyethylene oleate; alkylamines such as polyoxyethylene laurylamino ether, polyoxyethylene stearylamino ether, polyoxyethylene oleylamino ether, polyoxyethylene soybean aminoether and polyoxyethylene beef tallow aminoether; alkylamides such as polyoxyethylene lauric amide, polyoxyethylene stearic amide and polyoxyethyleneoleic amide; vegetable oil ethers such as polyoxyethylene castor oil ether and polyoxyethylene rapeseed oil ether; alkanol amides such as lauric acid diethanol amide, stearic acid diethanol amide and oleic acid diethanol amide; and sorbitan ester ethers such as polyoxyethylene sorbitan monolaurate, polyoxyethylene sorbitan monopalmitate, polyoxyethylene sorbitan monostearate and polyoxyethylene sorbitan monooleate. Of the above, polyoxyethylene stearates are preferred. Additional suitable surfactants can be found by reference to a standard text on surfactants such as those described in Ash and Ash, Encyclopedia of Surfactants, Chemical Publishing Company, New York, 1985; McCutcheon's Emulsifiers and Detergents, North American Edition, McCutcheon Publishing Company, Glen Rock, N.J., 2000; and Remington: The Science and Practice of Pharmacy, Nineteenth Edition, Vol. 1 at p. 251 coupled with the use of the procedures set forth above. [0034] The secondary surfactant is one or more surfactants meeting the above criteria and in addition, enables the emulsion to withstand autoclaving without significant degradation of the emulsion. The secondary surfactant desirably has a relatively small lipophilic group and a long chain hydrophilic group with a minimum of 6 repeating hydrophilic groups. More preferably, the secondary surfactant has an HLB of 9 or more, and most preferably, an HLB ranging between 12 and 20, a hydrophilic group of at least 9 repeating hydrophilic groups, most preferably at least 9 ethylene oxide groups or isopropyl oxide groups, and a relatively small lipophilic group that is the same or similar in structure to the lipophilic group of the primary surfactant. Exemplary nonionic surfactants include, but are not limited to the Octoxynol-n series of the formula C 8 H 17 C 6 H 4 (OCH 2 CH 2 )OH n where n is between 5 and 70 and preferably between 30 and 50, the nonoxynol-n series of the formula C 9 H 19 C 6 H 4 (OCH 2 CH 2 )OH p where p is between 5 and 40 and preferably between 15 and 30, and polyoxyethylene C 12-22 alkyl ethers such as polyoxyethylene lauryl ether, polyoxyethylene cetyl ether, polyoxyethylene stearyl ether or polyoxyethylene oleyl ether. Most preferred secondary surfactants are the Octoxynol series of surfactants having between 30 and 50 ethylene oxide groups. Numerous other surfactants having an HLB value of greater than about 9 and meeting the above criteria are listed in Ash and Ash, McCutcheons, and Remington, supra. [0035] The concentration of the surfactant combination used to form the emulsion may vary within wide limits. A treatment composition containing the surfactant combination in an amount as low as 0.01 weight percent of the total composition provides some benefit. A concentration of surfactant combination varying between 0.05 to 5.0 percent of the total composition is a clinically practical concentration range for purposes of the invention provided that the surfactant does not cause patient discomfort when used at the higher concentrations. Most preferably, the concentration of the combination varies between about 0.25 and 2.5 percent by weight of the total composition. It should be understood that with many surfactants, as concentration increases, the likelihood of physical discomfort—i.e., stinging, of the emulsion on the eye increases. Thus, if significant stinging occurs when the emulsion is applied to the ocular surface, it is likely that the concentration of surfactant is too high. [0036] The ratio of the primary surfactant to the secondary surfactant may vary within relatively broad limits—for example, between 0.2 to 1.0 to 1.0 to 0.2 primary to secondary surfactant. A more preferred range varies between 0.5 to 1.0 and 1.0 to 0.5. Most preferably, the primary surfactant is used in slightly larger concentration than the secondary surfactant and the most preferred ratio varies between 1.0 to about 0.75. [0037] The emulsions of the invention comprise an oil in water emulsion. The oil used to form the emulsion may be derived from animals, plants, nuts, petroleum, etc. Those derived from animals, plant seeds, and nuts are similar to fats and are primarily glycerides or fatty acids and consequently, contain a significant number of acid and/or ester groups rendering the same polar and lesser preferred for purposes of the invention. Alternatively, oils derived from petroleum are usually aliphatic or aromatic hydrocarbons that are essentially free of polar substitution and therefore suitable for purposes of the present invention provided the oil is refined so as to be compatible with human tissue such as the ocular surface. Preferably, the oil is a linear hydrocarbon oil having from 10 to 50 carbon atoms and more preferably, the oil is a saturated n-alkane or isoalkane hydrocarbon having from 14 to 26 carbon atoms. Unsaturated alkene hydrocarbons may be used but are less chemically stable. Aromatic oils are lesser preferred because it is known that aromatic compounds are for the most part unsuitable for application to the ocular surface. Mineral oil is the most preferred oil for purposes of this invention. [0038] The oil component within the emulsion may vary within reasonable limits provided the amount of oil retained on the eye following its application to the eye is within controlled volumes and does not exceed 25 μl, more preferably varies between about 1 and 10 μl and most preferably varies between about 1 and 5 μl. If the amount of oil added to the eye is in excess of 25 μl, the oil layer on the surface of the eye may be of excessive thickness and resulting in prolonged blurring. A treatment composition containing the oil in a concentration of at least 0.1 percent by weight of the total composition provides some benefits. A preferred concentration for the oil is at least 1.0 percent of the weight of the treatment composition. Preferably, the oil content of the treatment solution varies between about 2.5 and 12.5 percent by weight of the composition. [0039] Other additives may be present in the treatment composition. Such materials include minor amounts of neutral lipids and oils such as one or more triglycerides, cholesterol esters, the natural waxes and cholesterol; high molecular weight isoprenoids; stabilizers, additional surfactants; preservatives; pH adjusters to provide a composition preferably having a pH between about 6.5 and 7.8 and most preferably, between about 7.2 and 7.5; salt, glycerol or sugar in sufficient concentration to form an isotonic or mildly hypotonic composition; etc., all as would be obvious to those skilled in the art. [0040] Another useful class of additives comprises medications. As a consequence of the long term stability of the oil film formed over the surface of the eye using the emulsion compositions of the invention, prolonged and improved delivery of the medication to the eye results due to increased contact time of the medication on the eye. Medications suitable for delivery to the eye using the film forming compositions of the invention are those soluble in either the aqueous or oil phase of the composition though it is preferable that the medication be soluble in the oil phase. Illustrative medications include antibiotics, antiviral agents, anti-inflammatory agents and antiglaucoma agents such as illustrated in part in published European Patent Application No. 0 092 453 published Oct. 26, 1983, sections 5.3.1 and 5.3.2, incorporated herein by reference. [0041] Any additional additives are added to the emulsion are added prior to formation of the emulsion using simple mixing techniques. The concentration of the additive is dependent upon the specific additive, and preferably, total additive content in addition to the surfactant and the oil are at a maximum concentration level whereby the total weight of the organics in the oil phase does not exceed 15 percent of the total weight of the emulsion. [0042] In accordance with the invention, the emulsions may be made in accordance with standard procedures. Desirably, a commercial homogenizer is used to form the emulsion as equipment of this nature enhances the stability of the emulsion during transportation and storage. The use of commercial homogenizers for the formation of emulsions is within the skill of the art. [0043] The emulsions of the invention are also desirably used with subjects requiring ocular prostheses. In this instance, the treatment composition enhances the tear film layer and lubricates the boundary between the prosthesis and the ocular surface. When used with an ocular prosthesis, the treatment composition may be applied to the inner or both the inner and outer surfaces of the prostheses prior to insertion of the same into the eye. Regardless of how added, the amount available to form the oil layer should be within the limits set forth above. [0044] The invention will be better understood by reference to the examples that follow. In the examples, the thickness of the lipid layer of a tear film formed over the ocular surface is evaluated by projecting a light source onto the ocular surface while viewing the reflected images from the light source on a video screen. The light source and its location is one that illuminates a surface area on the ocular surface of approximately 10 mm 2 . Interference patterns are formed, the color(s) of which are indicative of the thickness of the oil layer. The color of the waves is correlated with a protocol of known film thickness. In this way, the tear film can be evaluated over a period of real time and rated in accordance with the following scale: [0000] Rating Film Characteristics Quality A Colored waves - particularly greens Excellent and blues. Waves extend from lower to above the lower pupillary border. Film thickness is excess of 170 nm. B Colored waves - reds, browns, yellows, Good but no blues. Waves extend from lower lid to above the pupillary border. Film thickness of approximately 90 nm. C Colored waves - only yellow is present. Good Waves extend form lower lid to lower pupillary border. Film thickness of approximately 90 nm. D Waves visible but no color present Fair or no color other than grayish white. Waves extend from lower lid to lower pupillary border. Film thickness of less than 90 nm. E No waves and no color. An absence of Poor any observable tear film movement. Film thickness of less than 70 nm. [0045] Further details pertaining to experimental procedure can be found in the above referenced U.S. Pat. Nos. 5,278,151; 5,294,607; and 5,578,586. [0046] The data presented in the examples was obtained using individuals with baseline lipid layers of C rating or less. The data illustrates the resultant change in lipid layer characteristics from the baseline finding to the finding for lipid characteristics after the application of a standard eye drop of the test formulation to the eye. A desirable result is for improvement in lipid layer characteristics, evidenced by an increase in the alphabetical grade, with A being the most desirable, and F being the least desirable. The evaluations were performed 5 minutes after the instillation of the test formulations. [0047] The first two examples illustrate that emulsions may be formed using surfactants having properties and HLBs suggesting suitability for formation of stable oil in water emulsions, but illustrate that the emulsions so formed are unstable at autoclaving temperatures and therefore unsuitable for dry eye treatment. Example 3 illustrates that a surfactant that might be suitable for formation of an emulsion having properties meeting the objectives of this invention is unsuitable as it causes discomfort to the patient when added to the eye. Examples 4 and 5 illustrate that primary surfactants unsuitable for formation of a dry eye emulsion can be made functional when used in combination with the secondary surfactants of the invention regardless of the arithmetic HLB. Example 1 [0048] This example illustrates that various surfactant combinations may be used that meet certain of the guidelines set forth above, especially those relating to HLB, but still fail to provide a meta stable emulsion able to withstand the elevated temperatures required for sterilization of the formulation. [0049] A mixture of Myrj-52, a polyoxyethylene (40) stearate, and glycerol monostearate (GMS), were used as primary surfactants to form a dry eye treatment emulsion with mineral oil as the oil phase. Myrj-52 has a high HLB value (15-16.9) and is water-soluble. Glycerol monostearate (GMS NF) has a low HLB value (3-5) and is therefore oil soluble suggesting that the combination should form a suitable oil in water emulsion. These surfactants have an identical lipophilic group (stearate) but different hydrophilic groups, and thus will have different physical behavior in terms of partitioning into the oil or water phases as suggested by the difference in the HLB value of the two surfactants. [0050] This combination of surfactants was evaluated by formation of 11 emulsions utilizing 5.5% (±0.3%) of Drakeol-35, a commercially available mineral oil at two different total surfactant concentrations −0.15% and 0.30%. The non-polar oil phase of Drakeol-35 mineral oil and the aqueous phase of 0.67% NaCl and 0.05% of anhydrous Na 2 HPO 4 were common to all 11 of these formulations; the pH was adjusted with diluted HCl as required. The relative concentration of the two individual surfactants was varied to evaluate the effect of the average HLB on emulsion quality as shown in Table 1. Emulsification of the 11 formulations was performed using a commercial homogenizer (PRO250) from Proscientific, Inc., with a ¾ horsepower motor which drove a 30 mm rotor-stator generator, by combining all of the reactants into one vessel and raising the temperature to approximately 90° C. Table 1 provides the formulations and the compositions (in grams) for the 11 test formulations utilizing the Myrj-52 and GMS primary surfactant systems. [0000] TABLE 1 gm Myrj-52 gm GMS gm Drakeol 35 Surfactant Content 0.094 0.065 5.700 0.15% 0.102 0.054 5.334 0.15% 0.112 0.042 5.508 0.15% 0.127 0.037 5.404 0.15% 0.130 0.030 5.284 0.15% 0.139 0.020 5.284 0.15% 0.298 0.017 5.309 0.30% 0.279 0.042 5.300 0.30% 0.258 0.065 5.283 0.30% 0.233 0.091 5.308 0.31% 0.207 0.108 5.303 0.30% Footnotes 1. Myrj-52 is polyoxyethylene (40) stearate 2. Glyceryl monostearate (GMS) is the glycerol ester of stearic acid 3. Drakeol refers to a series of NF mineral oils available from Penreco Co. of Butler, PA. The numeral following the letters represents the average molecular weight of the oil, and is an indication of the viscosity of the fluid. Results [0051] The formulations of Table 1 produced emulsions which all met the first criterion of separation when at rest for several minutes. They also appeared to meet the second criterion of the emulsion returning to its original dispersed form after simple agitation. However, after periods of 60 minutes to 1 day, some of the oil phase in the formulations with the higher HLB values evidenced significant oil breakout, where the individual oil droplets were broken, resulting in the formation of an oil film on the surface. [0052] The formulations with the lower HLB values provided emulsions that did not evidence the oil film after similar periods of time. However, when agitated with mechanical shaking to simulate transportation effects, the oil film was visible on the surface within a time period that precluded a commercially viable product. [0053] In general, an increasing value of the calculated HLB produced poorer quality emulsions upon standing or upon agitation. However, no formulation in Table 1 was found to be adequate because of the degradation of the individual oil droplets and the subsequent formation of a surface oil layer. Further, microscopic studies and photographs of these formulae taken before and after both autoclaving and shaking demonstrated that the oil droplets were degraded either by being subjected to autoclaving at 121° C. or by shaking on the Platform Rocker Shaker by Vari-Mix for less than one day. [0054] The above example illustrates that for the materials of this example, HLB values alone proved to be an unreliable parameter of complex formulation issues in the development of the intended dry eye treatment composition. While HLB values are generally useful as a formulation development guide, it was obvious that further considerations are required for the development of a dry eye treatment composition suitable for purposes of this invention. [0055] Resolution of the above described problem would likely require significantly higher surfactant concentrations but it was known that increasing the concentration of Myrj-52% would result in ocular discomfort—i.e. significant stinging. For this reason, the formulation was not evaluated clinically. Example 2 [0056] Polysorbate-80 (trade name Tween-80), a stearyl ether of a polysorbate, was evaluated as a sole primary surfactant for forming a dry eye treatment emulsion utilizing a total concentration of 7.0% of a mineral oil mixture of Draekol-15 and Draekol-35. In this study, Polysorbate-80 concentrations of 0.2%, 1.0%, and 1.5% were utilized in the base formula as displayed in Table 2 as formulae C1 through C6. These formulations were prepared with and without disodium EDTA. Emulsification of the 6 formulations was carried out with a commercial homogenizer (PRO250) from Proscientific, Inc., using a ¾ horsepower motor that drove a 30 mm rotor-stator generator, by combining all of the reactants into one vessel and raising the temperature to approximately 90° C. Table 2 provides the formulations and the compositions (in grams) for the 6 test formulations utilizing Tween-80 as the sole surfactant, [0000] TABLE 2 formulae C1 C1 C3 C3 C5 C6 D-15 2.07 2.07 2.02 2.02 2.02 2.03 D-35 5.17 5.17 5.04 5.04 5.05 5.07 Polysorbate-80 0.20 0.20 1.00 1.00 1.50 1.51 NaCl 0.67 0.67 0.67 0.67 0.67 0.67 Na 2 HPO 4 (anh.) 0.05 0.05 0.05 0.05 0.05 0.05 NaH 2 PO 4 •2H 2 O 0.02 0.02 0.02 0.02 0.02 0.03 EDTA 0.02 0.02 0.02 0.02 0.02 0.00 Water 100 100 100 100 100 100 Footnotes 1. Drakeol refers to a series of NF mineral oils available from Penreco Co. of Butler, PA. The numeral following the letters represents the average molecular weight of the oil, and is an indication of the viscosity of the fluid 2. Polysorbate-80 is a stearyl ether of a polysorbate and is sold under the tradename Tween-80 by ICI (now known as Uniqema, New Castle, DE) in Wilmington, DE. Results [0057] The six emulsions prepared using Polysorbate-80 as a sole primary surfactant were found to meet the first criterion of providing appropriate separation of the oil and aqueous phases upon resting in the container for several minutes. The second criterion of reconstitution by simple product agitation was also met. However, these formulae failed to meet the third criterion because all were destabilized (as evidenced by surface oil film formation) when agitated for short periods of time on the laboratory shaker or when autoclaved for 30 minutes at 121° C. Increased Polysorbate-80 content in the formulae to 1.50 decreased oil droplet instability, but still failed to meet the third criterion. The failure to meet the third criterion was confirmed in post autoclaved and post shaken samples, in that that the surface oil film formed by droplet coalescence prevented the reconstitution of the emulsion by simple shaking as required by the second criterion. [0058] These formulae thus did not meet the requirements for the desired eye treatment solution. It was found that when using only Polysorbate-80 as a sole primary surfactant, a higher concentration of Polysorbate-80 was required with the following undesirable results: (1) the emulsion was not autoclave stable, (2) the higher concentration of the Polysorbate-80 led to stinging on the eye and (3) the higher concentration of the Polysorbate-80 degraded the performance of the meta stable emulsion on the eye. Thus, thee performance of these samples failed by a wide margin to meet the criterion of maintaining the original emulsion characteristics after agitation for a period of time. Therefore, the use of Polysorbate-80 as a sole surfactant for the non-polar oil formula was judged inadequate and was not evaluated clinically. Example 3 [0059] Though this example does not illustrate the formation of an emulsion, it does illustrate the basis for rejection of a primary surfactant that would otherwise appear to be suitable for formation of an oil in water emulsion. [0060] The example describes the evaluation of Span 20, a highly viscous water insoluble sorbitan monolaurate, for suitability as a surfactant system for use on the eye. The HLB of 8.6 of this surfactant, and its reported use in ophthalmic products suggested that it would be a suitable surfactant for formation of an oil in water emulsion using mineral oil. [0061] Five concentrations of Span 20, from 0.05% to 1.00% WV were prepared in a buffered normal saline vehicle for evaluation of comfort on the eye. The vehicle used for all formulations was Unisol 4, a buffered saline solution marketed by Alcon Laboratories, Fort Worth, Tex. Unisol 4 had previously been studied and found to be the most comfortable normal saline product for use on the ocular surface. It was therefore used as a vehicle for the test formulations, and was also used as a control. [0062] A drop of each of the 5 test formulations was placed on to the ocular surfaces of each subject, utilizing a 15 ml dropper container that delivered a drop of 40 μl to the ocular surface. The subject was asked to describe the sensation as one of: pleasant, neutral, slight sting, moderate sting, or severe sting. Results [0063] The results obtained with 6 subjects, all of whom evaluated each of the test formulation on two different days are summarized in Table 3. [0000] TABLE 3 Formulation Results S477 (Control, Unisol 4) Pleasant to neutral S478 (0.05% Span 20) Neutral to slight sting S479 (0.10% Span 20) Neutral to slight sting S480 (0.20% Span 20) Slight to moderate sting S481 (0.40% Span 20) Moderate sting S482 (1.00% Span 20) Moderate to severe sting Footnote 1. Span 20 is a highly viscous water insoluble sorbitan monolaurate sold under the tradename Span 20 by ICI (now known as Uniqema, New Castle, DE) in Wilmington, DE. [0064] In view of the sting when applied to the eye and ocular surfaces, even in concentrations ≦0.10%, the use of Span 20 was rejected as a suitable surfactant for an eye treatment composition. Example 4 [0065] This example illustrates that Polysorbate-80 found unsatisfactory for purposes of this invention in Example 2 can be made suitable by combination with a secondary surfactant, in this case Octoxynol-40. [0066] The example determines the optimum ratio of a mixed surfactant system comprising Polysorbate-80 as a primary and Octoxynol-40 as a secondary surfactant to maximize the stability of the oil/water interface in the ocular emulsion systems. Formulae utilizing 2.4% Drakeol-15 and 4.8% Drakeol-35 with both Octoxynol-40 and Polysorbate-80 were evaluated at different concentrations. Octoxynol-40 has an HLB of 19 and Polysorbate-80 has an HLB of 15. The two surfactants combined will yield an arithmetic HLB above the HLB believed suitable for formation of a stable oil in water emulsion. The samples were made both with and without disodium EDTA, keeping the concentrations of the other additives at the levels required for eye treatment compositions. The emulsification of the four formulations prepared was carried out with a commercial homogenizer (PRO250) from Proscientific, Inc., using a ¾ horsepower motor which drove a 30 mm rotor-stator generator, by combining all of the reactants into one vessel and raising the temperature to approximately 60° C. Table 4 sets forth the formulations and the compositions (in grams) for the 4 test formulations utilizing the Octoxynol-40 and Polysorbate-80 surfactant systems. [0000] TABLE 4 Formulae D1 (AD1) D2 (AD2) D3 (AD3) D4 (AD4) D-15 2.42 2.51 2.41 2.31 D-35 4.83 4.85 4.86 4.84 Octoxynol-40 1.12 1.12 0.63 0.63 Polysorbate-80 0.38 0.38 0.90 0.91 NaCl 0.67 0.67 0.67 0.67 Na 2 HPO 4 (anh.) 0.05 0.05 0.05 0.05 NaH 2 PO 4 •2H 2 O 0.03 0.02 0.02 0.03 EDTA 0 0.02 0.02 0 Water 100 100 100 100 Footnotes 1. Drakeol refers to a series of NF mineral oils available from Penreco Co. of Butler, PA. The numeral following the letters represents the average molecular weight of the oil, and is an indication of the viscosity of the fluid. 2. Octoxynol-40 is polyethylene glycol (40) p-isooctylphenyl ether sold under the tradename Synperonic OP-40 by ICI (now known as Uniqema, New Castle, DE) in Wilmington, DE. 3. Polysorbate-80 is a stearyl ether of a polysorbate sold under the tradename Tween-80 by ICI (now known as Uniqema, New Castle, DE) in Wilmington, DE. Results [0067] All 4 formulations met the three pre-clinical criteria of proper separation, reconstitution of the emulsion by gentle shaking and maintaining the original emulsion characteristics when agitated for a period of time of at least 72 hours on the laboratory Platform Rocker Shaker by Vari-Mix, or when autoclaved for 30 minutes at 121° C. [0068] Microscopic studies and photographs of the samples-before and after autoclaving and shaking demonstrated that the oil droplets were not degraded by being subjected to autoclaving at 121° C., or by shaking on the Platform Rocker Shaker by Vari-Mix for 288 hours. [0069] Since this formulation of non-polar oil with the surfactant system of Octoxynol-40 and Polysorbate-80 satisfied the pre-clinical criteria, formulae D1 and D3 were evaluated clinically and were found to adequately augment and restore the lipid layer thickness. The results of the clinical evaluations of formulations D1 and D4 are given in the following table: [0000] Rating before Rating after Patient Number Formulation treatment treatment 1 D1 C A 2 D1 D B 3 D1 D A 4 D1 D B 5 D1 C A 6 D1 C B 7 D1 C A 8 D4 D B 9 D4 D C 10 D4 D B 11 D4 D A 12 D4 C A 13 D4 C A 14 D4 C A [0070] The clinical evaluations of patient numbers 1 to 14 indicated that the surfactant system of Octoxynol-40 and Polysorbate-80 were efficacious in forming and restoring a lipid layer of improved characteristics and that both formulations were essentially equally effective. The subjective sensation realized with both formulations was evaluated as comfortable, and without any form of adverse sensation. The studies of example 3 indicated that the addition of the surfactant Octoxynol-40 as a second surfactant in combination with Polysorbate-80 satisfied the criteria for a dry eye treatment composition. Example 5 [0071] This study was directed to determining the optimal concentrations of the mixed primary and secondary surfactant system of Polysorbate-80 (Tween-80) and Octoxynol-40 for an optimal dry eye treatment formulation. [0072] Six formulations were prepared. Each was formed using a commercial homogenizer (PRO250) from Proscientific, Inc., using a ¾ horsepower motor which drove a 30 mm rotor-stator generator, by combining all of the reactants into one vessel and raising the temperature to approximately 60° C. Table 5 illustrates 6 formulations, where the concentration of Polysorbate-80 is held constant at 0.40%, while the concentration of Octoxynol-40 is varied at 0.30%, 0.60% and 1.20% surfactant levels. The 0.40% Polysorbate-80 concentration was chosen from the result of prior experiments that established that this concentration, when used with higher levels of Octoxynol-40, met the pre-clinical requirements, while lower concentrations of Polysorbate-80 had resulted in minimal but detectable oil droplet degradation after being subjected to autoclaving at 121° C., or by shaking on the Platform Rocker Shaker by Vari-Mix for periods of time from 3 days to 288 hours. It was also considered desirable to utilize the lowest concentrations of both Polysorbate-80 and Octoxynol-40 for a dry eye treatment composition, while meeting the previously described requirements, since the sensitivity of the ocular surface cells and the immune response of the eye can be expected to increase with increasing concentration to any compound placed on the eye. These phenomena may also be exacerbated in dry eye states, since the surface epithelial cells are frequently compromised by the desiccation and lack of lubrication accompanying dry eye states. Table 5 provides the formulations and the compositions (in grams) for the 6 test formulations: [0000] TABLE 5 Formulae F1 F2 G1 G2 H1 H2 Polysorbate-80 0.41 0.41 0.40 0.40 0.40 0.40 Octoxynol-40 1.21 1.21 0.60 0.60 0.30 0.30 EDTA•2H 2 0 0.02 0.02 0.02 0.02 0.02 0.02 NaHPO 4 0.05 0.05 0.05 0.05 0.05 0.05 NaH 2 PO 4 •2H 2 0 0.02 0.02 0.02 0.02 0.02 0.02 NaCl 0.67 0.67 0.67 0.67 0.67 0.67 H 2 0 100 100 100 100 100 100 Drakeol-15 2.17 2.16 2.14 2.15 2.12 2.13 Drakeol-35 5.42 5.41 5.36 5.38 5.31 5.33 P/O 0.5 0.5 1.0 1.0 2.0 2.0 rpm 4050 10000 4050 10000 4050 10000 Footnotes 1. Polysorbate-80 is a stearyl ether of a polysorbate sold under the tradename Tween-80 by ICI (now known as Uniqema, New Castle, DE) in Wilmington, DE. 2. Octoxynol-40 is polyethylene glycol (40) p-isooctylphenyl ether sold under the tradename Synperonic OP-40 by ICI (now known as Uniqema, New Castle, DE) in Wilmington, DE. 3. Drakeol refers to a series of NF mineral oils available from Penreco Co. of Butler, PA. The numeral following the letters represents the average molecular weight of the oil, and is an indication of the viscosity. 4. P/O is the molar ratio of the Polysorbate-80 to Octoxynol-40 components. 5. RPM is the speed of the homogenization unit. [0073] The lowest concentration of 0.30 grams Octoxynol-40 with 0.40 grams Polysorbate-80 was found to optimally satisfy the 3 criteria. The low surfactant concentration is desirable due to the possibility of irritancy of ocular tissue, where the sensitivity to a compound placed on the eye can be expected to increase with increasing concentration.

This invention relates to an emulsion composition for the formation of an artificial tear film over the ocular surface of the eye capable of providing mechanical lubrication for the ocular surface while reducing evaporation of fluid therefrom. The emulsion is desirably in the form of a meta stable emulsion and is characterized by the use of a surfactant comprising a combination of a primary and secondary surfactant where the primary surfactant permits formation of the emulsion and the secondary surfactant permits autoclaving of the surfactant. The invention also includes a method for the formation of such an emulsion.

RELATED APPLICATION DATA This application is a continuation-in-part of copending provisional application No. 60/000,496 filed Jun. 26, 1995 and a continuation of International Application No. PCT/US96/10899, filed Jun. 26, 1996. FIELD OF THE INVENTION The herein described invention relates generally to a cushioning conversion machine and method for converting sheet-like stock material into a cushioning product. BACKGROUND OF THE INVENTION In the process of shipping an item from one location to another, a protective packaging material is typically placed in the shipping case, or box, to fill any voids and/or to cushion the item during the shipping process. Some conventional protective packaging materials are plastic foam peanuts and plastic bubble pack. While these conventional plastic materials seem to adequately perform as cushioning products, they are not without disadvantages. Perhaps the most serious drawback of plastic bubble wrap and/or plastic foam peanuts is their effect on our environment. Quite simply, these plastic packaging materials are not biodegradable and thus they cannot avoid further multiplying our planet's already critical waste disposal problems. The non-biodegradability of these packaging materials has become increasingly important in light of many industries adopting more progressive policies in terms of environmental responsibility. The foregoing and other disadvantages of conventional plastic packaging materials have made paper protective packaging material a very popular alterative. Paper is biodegradable, recyclable and composed of a renewable resource, making it an environmentally responsible choice for conscientious industries. While paper in sheet form could possibly be used as a protective packaging material, it is usually preferable to convert the sheets of paper into a relatively low density pad-like cushioning dunnage product. Cushioning conversion machines in use today have included a forming device and a feeding device which coordinate to convert a continuous web of sheet-like stock material (either single-ply or multi-ply) into a three dimensional cushioning product, or pad. The forming device is used to fold, or roll, the lateral edges of the sheet-like stock material inward on itself to form a strip having a width substantially less than the width of the stock material. The feeding device advances the stock material through the forming device and it may also function as a crumpling device and a connecting (or assembling) device. The cushioning conversion machine may also include a ply-separating device for separating the plies of the web before passing through the former, and usually a severing assembly; for example, a cutting assembly for cutting the strip into sections of desired length. European Patent Application No. 94440027.4 discloses a cushioning conversion machine wherein the feeding device comprises input and output pairs of wheels or rollers which operate at different speeds to effect, along with feeding of two plies of paper, crumpling and assembling of the paper plies to form a connected strip of dunnage. The cushioning conversion art would benefit from improvements in the machine shown in such application, and such improvements may have applicability to other cushioning conversion machines as well. SUMMARY OF THE INVENTION The present invention provides an improved cushioning conversion machine and related methodology characterized by one or more features including, inter alia, a feeding/connecting assembly which enables an operator to easily vary a characteristic, for example, the density, of the cushioning product; a feeding/connecting assembly wherein input and/or output wheels or rollers thereof are made at least in part of an elastomeric or other friction enhancing material, which reduces the cost and complexity of the input and output rollers; a manual reversing mechanism that is useful, for example, for clearing paper jams; a modular arrangement of a forming assembly and feeding/connecting assembly in separate units that may be positioned remotely from one another, as may be desired for more efficient utilization of floor space; a layering device which provides for doubling of the layers of sheet material in the converted cushioning product; a turner bar which enables alternative positioning a stock supply roll; and a volume expanding arrangement cooperative with the feeding/connecting assembly for reducing the density of the cushioning product and increasing product yield. The features of the invention may be individually or collectively used in cushioning conversion machines of various types. These and other aspects of the invention are hereinafter summarized and more fully described below. According to one aspect of the invention, a cushioning conversion machine, for making a cushioning product by converting an essentially two-dimensional web of sheet-like stock material of at least one ply into a three-dimensional cushioning product, generally comprises a housing through which the stock material passes along a path; and a feeding/connecting assembly which advances the stock material from a source thereof along said path, crumples the stock material, and connects the crumpled stock material to produce a strip of cushioning. The feeding/connecting assembly includes upstream and downstream components disposed along the path of the stock material through the housing, at least the upstream component being driven to advance the stock material toward the downstream component at a rate faster than the sheet-like stock material can pass from the downstream component to effect crumpling of the stock material therebetween to form a strip of cushioning. Additionally, at least one of the upstream and downstream components includes opposed members between which the stock material is passed and pinched by the opposed members with a pinch pressure; and a tension control mechanism is provided for adjusting the amount of pinch pressure applied by the opposed members to the stock material. In one embodiment of the invention, the tension control mechanism includes an accessible control member outside the housing for enabling easy operator adjustment of the pinch pressure, whereby a characteristic of the strip of cushioning can be varied on demand. In another embodiment, the upstream and downstream components each include opposed members between which the stock material is passed and pinched by the opposed members with a pinch pressure; and a tension control mechanism is provided for adjusting the amount of pinch pressure applied to the stock material by the opposed members of the downstream component independently of the pinch pressure applied to the stock material by the opposed members of the upstream component, whereby a characteristic of the strip of cushioning can be varied. According to another aspect of the invention, a cushioning conversion machine again generally comprises a housing through which the stock material passes along a path; and a feeding/connecting assembly which advances the stock material from a source thereof along the path, crumples the stock material, and connects the crumpled stock material to produce a strip of cushioning. The feeding/connecting assembly includes upstream and downstream feeding components disposed along the path of the stock material through the housing, the upstream feeding component being driven to advance the stock material toward the downstream component at a rate faster than the sheet-like stock material can pass from the downstream component to effect crumpling of the stock material therebetween to form the strip of cushioning. An adjustable speed control mechanism is provided for varying the ratio of the feeding speeds of the upstream and downstream feeding components, whereby a characteristic of the strip of cushioning can be varied. In a preferred embodiment, the adjustable speed control mechanism can include, for example, a variable speed drive device (such as a variable pitch pulley systems for one of the upstream and downstream components, a quick change gear set, or a variable speed control for at least one of respective drive motors for the upstream and downstream components. Preferably, a control member is provided outside the housing for enabling easy operator adjustment of the speed ratio, whereby a characteristic of the strip of cushioning can be varied on demand. According to a further aspect of the invention, a cushioning conversion machine again generally comprises a housing through which the stock material passes along a path; and a feeding/connecting assembly which advances the stock material from a source thereof along the path, crumples the stock material, and connects the crumpled stock material to produce a strip of cushioning. The feeding/connecting assembly includes upstream and downstream components disposed along the path of the stock material through the housing, at least the upstream component being driven to advance the stock material toward the downstream component at a rate faster than the sheet-like stock material can pass from the downstream component to effect crumpling of the stock material therebetween to form a strip of cushioning. Also provided is a stretching component downstream of the downstream component that is operative to advance the strip of cushioning at a rate faster than the rate at which the stock material passes from the downstream component to effect longitudinal stretching of the strip of cushioning. According to yet another aspect of the invention, a cushioning conversion machine again generally comprises a housing through which the stock material passes along a path; and a feeding/connecting assembly which advances the stock material from a source thereof along the path, crumples the stock material, and connects the crumpled stock material to produce a strip of cushioning. The feeding/connecting assembly includes upstream and downstream components disposed along the path of the stock material through the housing, at least the upstream component being driven to advance the stock material toward the downstream component at a rate faster than the sheet-like stock material can pass from the downstream component to effect crumpling of the stock material therebetween to form a strip of cushioning. At least one of the upstream and downstream components includes opposed members between which the stock material is passed and pinched by the opposed members with a pinch pressure; and at least one of the opposed members is at least partially made of an elastomeric material at a surface thereof engageable with the stock material. According to a still further aspect of the invention, a cushioning conversion machine generally comprises a housing through which the stock material passes along a path; and a feeding/connecting assembly which advances the stock material from a source thereof along the path, crumples the stock material, and connects the crumpled stock material to produce a strip of cushioning. The feeding/connecting assembly includes at least one rotatable member rotatable in a first direction for engaging and advancing the stock material along the path, a feed motor for driving the one rotatable member in the first direction, and a crank coupled to the rotatable member for enabling rotation of the one rotatable member in a second direction opposite the first direction. In a preferred embodiment the crank is coupled to the rotatable member by a one-way clutch. According to yet still another aspect of the invention, a cushioning conversion machine comprises first and second units having separate housings whereby the first and second units can be located at spaced apart locations. The first unit includes in the housing thereof a former for folding the sheet-like stock material to form flat folded stock material having a plurality of layers each joined at a longitudinally extending fold to at least one other layer. The second unit includes in the housing thereof an expanding device operative, as the flat folded stock material passes therethrough, to separate adjacent layers of the flat folded stock material from one another to form an expanded strip of stock material, and a feeding/connecting assembly which advances the stock material through the expanding device, crumples the expanded stock material passing from the expanding device, and connects the crumpled strip to produce a strip of cushioning. In a preferred embodiment, the units are used in combination with a table to form a packaging system, the table including a table top having a packaging surface. The first and second units may be both located beneath said packaging surface, and one may be supported atop the other. In alternative arrangement, the first unit may be located beneath the table top and the second unit may supported on the table top. According to another aspect of the invention, a cushioning conversion machine generally comprises a supply assembly for supplying the sheet-like stock material; and a conversion assembly which converts the sheet-like stock material received from the supply assembly into a three-dimensional strip of cushioning. The stock supply assembly includes a support for a supply of the stock material from which the stock material can be dispensed, and a layering device which effects folding of the stock material along a fold line parallel to the longitudinal axis of the stock material, thereby in effect doubling the number of layers of the stock material that are converted into a cushioning product. According to a further aspect of the invention, a cushioning conversion machine comprises a forming assembly through which the sheet-like stock material is advanced to form the stock material into a three-dimensional shape and a feeding/connecting assembly that advances and crumples the formed strip, and connects the crumpled formed strip to produce a strip of cushioning. The forming assembly includes a forming member and a converging chute cooperative with the forming member to cause inward rolling of the edges of the stock material to form lateral pillow-like portions of a formed strip, and the feeding/connecting assembly includes upstream and downstream components disposed along the path of the stock material through the machine, at least the upstream component being driven to advance the stock material toward the downstream component at a rate faster than the sheet-like stock material can pass from the downstream component to effect crumpling of the stock material therebetween to form a strip of cushioning. According to yet another aspect of the invention, a cushioning conversion machine comprises a feeding/connecting assembly which advances the stock material from a source thereof along a path through the machine, crumples the stock material, and connects the crumpled stock material to produce a strip of cushioning. The feeding/connecting assembly includes upstream and downstream feeding components disposed along the path of the stock material through the housing, the upstream feeding component being driven continuously to advance continuously the stock material toward the downstream feeding component during a cushioning formation operation, and the downstream feeding component being driven intermittently to advance periodically the stock material. Accordingly, when the downstream feeding component is not driven the stock material will be caused to crumple longitudinally between the upstream and downstream feeding components, and when driven the longitudinally crumpled stock material will be advanced by the downstream feeding component toward an exit end of the machine. According to a still further aspect of the invention, a method for making a cushioning product, by converting an essentially two-dimensional web of sheet-like stock material of at least one ply into a three-dimensional cushioning product, generally includes the steps of supplying the stock material, and using an upstream component of a feeding/connecting assembly to advance the stock material toward a downstream component of the feeding/connecting assembly at a rate faster than the stock material can pass from the downstream component to effect crumpling of the stock material therebetween to form the strip of cushioning, the upstream and downstream components including opposed members between which the stock material is passed and pinched by the opposed members with a pinch pressure. In one embodiment, the method includes the step of adjusting the amount of pinch pressure applied by the opposed members of the downstream component independently of the pinch pressure applied to the stock material by the opposed members of the upstream component to the stock material, whereby a characteristic of the strip of cushioning can be varied. In another embodiment, the method includes the step of varying the ratio of the feeding speeds of the upstream and downstream feeding components, whereby a characteristic of the strip of cushioning can be varied. The foregoing and other features of the invention are hereinafter fully described and particularly pointed out in the claims, the following description and the annexed drawings setting forth in detail certain illustrative embodiments of the invention, these being indicative, however, of but a few of the various ways in which the principles of the invention may be employed. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is a top plan view of a cushioning conversion machine according to the present invention, the machine including a housing, stock-supply assembly, a forming assembly, a feeding/connecting assembly, a severing assembly, and a post-severing assembly. FIG. 2 is a schematic side elevational view of the cushioning conversion machine 100. FIG. 3 is a sectional view of the feeding/connecting assembly of the machine 100 and relevant portions of the machine's housing. FIG. 3A is a fragmentary view of a gear of the feeding/connecting assembly and a relevant portion of the machine's housing. FIGS. 4A and 4B are edge and side views, respectively, of a component of the feeding/connecting assembly, namely a feed wheel. FIGS. 4C and 4D are edge and side views, respectively, of a component of the feeding/connecting assembly, namely a support wheel for the feed wheel. FIGS. 4E and 4F are edge and side views, respectively, of a component of feeding/connecting assembly, namely a compression wheel. FIGS. 4G and 4H are edge and side views, respectively, of a component of the feeding/connecting assembly, namely a support wheel for a compression wheel. FIG. 5A is an isolated plan view of the feeding/connecting assembly, along with relevant parts of the machine's frame or housing. FIG. 5B is a side view of the feeding/connecting assembly, as seen from the line 5B--5B in FIG. 5A. FIG. 5C is a sectional view of the feeding/connecting assembly, taken along line 5C--5C of FIG. 5A. FIGS. 6A and 6B are schematic side and plan views, respectively, of another cushioning conversion machine 100 according to the present invention., FIG. 6C is schematic side view of the forming assembly of the cushioning conversion machine. FIG. 7 is a side view of portions of a modified version of the feeding/connecting assembly of FIGS. 1-2. FIG. 8 is a side view of portions of a modified version of the feeding/connecting assembly of FIGS. 1-2. FIG. 9 is a sectional view taken along line 9--9 in FIG. 8. FIG. 10 is a schematic view of portions of a modified version of the feeding/connecting assembly of FIGS. 1-2. FIGS. 11A and 12 are schematic plan view of first and second modular units of another cushioning conversion machine according to the present invention. FIG. 11B is an end view of device of the first modular unit, namely an expanding device, the device being shown with flat-folded stock material expanded thereby. FIG. 11C is a side view of the expanding device of FIG. 11B, without the stock material. FIGS. 13-15 are side elevation view of three packaging systems according to the present invention which incorporates the cushioning conversion machine shown in FIGS. 11A and 12. FIG. 16 is a side elevation view of a packaging system according to the present invention which incorporates a modified version of the second modular unit shown in FIG. 12. FIG. 17 is a partial plan view of a modified version of the stock supply assembly of FIGS. 1-2. FIG. 18 is side elevation view of the modified version of the stock supply assembly of FIG. 17. FIG. 19A is a plan view of a modified version of the feeding/connecting assembly of FIGS. 1 and 2. FIG. 19B is a side elevation view of the feeding/connecting assembly of FIG. 19A. FIG. 19C is a cross-sectional view of the feeding/connecting assembly of FIG. 19A, the section being taken along line 19C--19C in FIG. 19A. FIG. 20 is a side elevation view of a modified version of the feeding/connecting assembly of FIGS. 1 and 2. FIG. 21 is an end elevation view of the feeding/connecting assembly of FIG. 20. FIG. 22 is a plan elevation view of a modified version of the feeding/connecting assembly of FIGS. 1 and 2. FIG. 23 is a cross sectional view of the feeding/connecting assembly of FIG. 22, the section being taken along fine 23--23 in FIG. 22. FIG. 24 is an end view of the feeding/connecting assembly of FIG. 22. DETAILED DESCRIPTION In FIGS. 1 and 2, a cushioning conversion machine 100 according to the present invention is shown. The machine 100 converts an essentially two-dimensional web of sheet-like stock material (the thickness thereof being negligible compared to the width and length thereof--thus the phrase "essentially two-dimensional) into a three-dimensional cushioning product of a desired length. The preferred stock material consists of plural plies or layers of biodegradable and recyclable sheet-like stock material such as 30 to 50 pound Kraft paper rolled onto a hollow cylindrical tube to form a roll R of the stock material. More preferably, the stock material consists of two plies of paper which are intermittently glued together with small drops of glue up the center of the paper plies, the glue drops being spaced approximately one foot apart. The preferred cushioning product has lateral accordion-like or pillow-like portions and is connected, or assembled, along a relatively thin central band separating the pillow-like portions. The cushioning conversion machine 100 includes a housing 102 having a base plate or wall 103, side plates or walls 104, a downstream end plate or wall 105, a top cover 106, and a downstream cover, or wall 107. The base, side, and end walls 103-105 collectively form the machine's frame structure. The top cover 106, together with the base, side and end walls 103-105, form an enclosure for the interior assemblies of the machine 100. (It should be noted that the terms "upstream" and "downstream" in the context of the present application correspond to the direction of flow of the stock material through the machine 100.) The walls 103-107 of the housing 102 are each generally planar and rectangular in shape. The upstream edges of the base wall 103 and sides walls 104 are turned in to form, along with a top bar 108, a rectangular border defining a centrally located, and relatively large, rectangular stock inlet opening. The rectangular border may be viewed as an upstream end plate or wall extending perpendicularly from the upstream edge of the base wall 103. The end plate 105 extends perpendicularly from a location near, but inward from, the downstream end of the base wall 103 and defines a dunnage outlet opening. The downstream cover wall 107 is attached to the downstream edges of the base wall 103, with the side walls 104 and a downstream portion of the top cover 106 forming a box-like enclosure for certain components of the machine 100. Preferably, the cover wall 107 may be selectively opened to provide access to these components. The downstream portion of the top cover preferably is fixedly secured in place while an upstream portion of the top cover may be in the form of a hinged door which may be opened to gain access to the interior of the housing and particularly the below mentioned forming assembly to facilitate loading of the stock material in a well known manner. The cushioning conversion machine 100 further includes a stock supply assembly 109, a forming assembly 110, a feeding/connecting assembly 111, a severing assembly 112, and a post-severing assembly 113. During the preferred conversion process, the stock supply assembly 109 supplies stock material to the forming assembly 110. The forming assembly 110 causes inward folding of lateral edge portions of the sheet-like stock material into an overlapping relationship. The feeding/connecting assembly 111 advances the stock material through the machine 100 and also crumples the folded over stock material to form a dunnage strip. As the dunnage strip travels downstream from the feeding/connecting assembly 111, the severing/aligning assembly 112 severs or cuts the dunnage strip into sections, or pads, of a desired length. The cut pads then travel through the post-severing assembly 113. The stock supply assembly 109 includes support brackets 114 which are laterally spaced apart and mounted to the upstream end of the machine's housing 102. The stock supply assembly 109 also includes first and second guide rollers 115 and 116 which are rotatably mounted between the support brackets 114, and a dancer roller 117 which is pivotally suspended from the support brackets 114 via swing arms 118. As paper is unwound from the stock or supply roll R, it travels around the dancer roller 117 so that the pull of the paper upward on the dancer roller 117, combined with the pull of gravity downward on the dancer roller and swing arms 118, helps maintain a uniform tension on the paper. The paper then travels over and under the two guide rollers 115 and 116 to guide the paper into the forming assembly 110. The forming assembly 110 consists of a central plate 119, a pair of fold-down rollers 120, with folding elements 121 and 122 forming a chute-like passage, or chute, for lateral edge portions of the stock material. The central plate 119 is mounted on a pedestal 123 attached to the base wall 103 and slopes slightly downwardly, and tapers inwardly, going from the upstream end to the downstream end of the central plate. The rollers 120 are mounted on a shaft 124a extending between the ends of a pair of swing arms 124b that are pivotally connected at their opposite ends to a support bar 124c extending between the side walls 104. The folding elements 121 and 122 are mounted, in a cantilever-like fashion, from a mounting plate 125. As the paper enters the forming assembly 110, the central portion of the paper (preferably about 1/3 of the paper width) will be positioned on the central plate 119 and its remaining lateral edge portions (preferably each about 1/3 the paper width) will be urged, or folded, downward by the rollers 120. As the paper contacts the folding elements 121 and 122, the folding elements will fold the lateral edge portions of the paper inward one over the other, whereby they will overlap in a folded arrangement. This overlapped paper, or strip, advances to the feeding/connecting assembly 111. The feeding/connecting assembly 111 includes a support structure 126, a wheel (or roller) network 127, a drive system 128, and a guide chute 129. The feeding/connecting components 126-129 feed the stock material, for example by pulling it from the stock supply assembly 109 and through the forming assembly 110. The feed/connecting assembly 111 longitudinally crumples the strip of stock material and then connects, or assembles, overlapped portions of stock material together to lock in a desired three-dimensional geometry of the resultant pad. With additional reference to FIGS. 3 and 5A-5C, the support structure 126 includes a pair of vertical side plates 130, and a horizontal cross bar 131. The downstream edges of the side plates 130 are coupled to the machine's housing 102, and more particularly to the end wall 105. The cross bar 131 extends between and is secured to the side plates 130. As best shown in FIGS. 3 and 5A-5C, the wheel network 127 includes a feed (or input) wheel 132, a support wheel 133 for the feed wheel 132, a compression (or output) wheel 134, a support wheel 135 for the compression wheel 134, and shafts 137-140 for each of the wheels 132-135, respectively. The lower wheels 132 and 134 are secured to the shafts 137 and 139, respectively, and the upper wheels 133 and 135 are rotatably mounted on their shafts 138 and 140, respectively. During operation of the feeding/connecting assembly 111, the lower shafts 137 and 139 are positively driven by the drive system 128 to rotate the lower wheels 132 and 134 which will in turn rotate the upper, or "idler", wheels 133 and 135. The lower shafts 137 and 139 extend between, and are rotatably journalled in the support side plates 130. (See FIGS. 3 and 5A-5C.) The upper shaft 140 extends between the side plates 130 and has its opposite ends positioned within a vertical guide slot 130a in the corresponding side plate 130. (See FIGS. 3 and 5A-5B.) The upper shaft 138 has opposite ends thereof terminating short of the side plates. A pair of laterally spaced apart shaft connectors 142 are connected between the upper shafts 138 and 140, and each shaft connector is attached, at about the middle thereof, to the lower end of a respective suspension pin or member 143. Each pin extends vertically though a respective guide opening in the cross bar 131 and carries thereon a compression spring 144 interposed between the cross bar and shaft connector. In this manner, the upper or "idler" wheels 133 and 135 will be resiliently biased towards the corresponding lower wheels 132 and 134, while being able to vertically "float" relative thereto during operation of the machine 100. As seen in FIGS. 4A-4D, the wheels 132 and 133 are both generally cylindrical in shape. The feed wheel 132 includes a middle portion 145 separating opposite axial end portions 146. The middle portion 145 is in the form of an annular groove which, for example, may have an approximately rectangular (as shown) or semi-circular cross section. The cylindrical periphery of the opposite axial end portions 146 is interrupted by flat faces 147. The flat faces 147 on one end portion 146 are staggered relative to the flat faces on the other end portion 146. In other words, the flat faces 147 on one axial end portion 146 are aligned with the "non-flat", or arcuate, knurled areas 148 on the other axial end portion 146. The support wheel 133 for the feed wheel 132 also includes a middle portion 149 separating opposite axial end portions 150. The middle portion 149 is in the form of a radially outwardly protruding annular rib which is preferably rounded at its radial outer side, while the end portions 150 have knurled radial outer surfaces. The radial outer surfaces of one or both of the wheels 132 and 133, or portions thereof, may be manufactured from an elastomeric material, such as rubber (neoprene or urethane) thereby reducing the cost and complexity of the wheels while still providing a high level of friction-enhancement for relatively slip free engagement with the stock material. As seen in FIGS. 4E-4H, the wheels 134 and 135 are also both generally cylindrical in shape. The compression wheel 134 includes a middle portion 151 separating opposite axial end portions 152. The middle portion 151 is radially relieved and has a smooth radial surface. The end portions 152 are ribbed to form rectangular, circumferentially spaced apart teeth. The support wheel 135 for the compression wheel 134 includes a continuous, knurled outer diameter surface. The radial outer surfaces of one or both of the wheels 134 and 135, or portions thereof, may again be manufactured from an elastomeric material such as rubber (neoprene or urethane) thereby reducing the cost and complexity of the wheels while still providing a high level of friction-enhancement for relatively slip free engagement with the stock material. As seen in FIG. 1, the drive system 128 for the feeding/connecting assembly 111 includes an electric motor 153, and motion-transmitting elements 154-159 (FIGS. 3, 3A and 5A). The motor 153 is mounted to the base plate 103 on one side of the forming assembly 110. The motion-transmitting elements transfer the rotational power of the motor 153 to the wheel network 127, or more particularly the lower shafts 137 and 139. As seen in FIGS. 3, 3A and 5A, the motion-transmitting elements include a drive chain 154 and sprockets 155 and 156. The sprocket 155 is secured to an output shaft 153a of a speed reducing gear box 153b driven by the motor 153 (See FIG. 1), and the sprocket 156 is secured to the compression wheel shaft 139. The drive chain 154 is trained around the sprockets 155 and 156 to rotate the compression wheel shaft 139. The motion transmitting elements 157-159 are gears forming a gear train between the compression wheel shaft 139 and the feed wheel shaft 137. The gear 157 is secured to the end of the compression wheel shaft 139 opposite the sprocket 156, the gear 158 is rotatably mounted to support side plate 130, and the gear 159 is secured to an adjacent end of the feed wheel shaft 137. In this manner, the feed wheel shaft 137 and the compression wheel shaft 139 will rotate in the same direction. However, the gears are selected so that the shaft 137 (and thus the feed wheel 132) is rotating at a faster feed rate than the shaft 139 (and thus the compression wheel 134). In the illustrated embodiment, the set speed ratio is on the order of about 1.7:1 to about 2.0:1. As seen in FIGS. 1 and 2, the guide chute 129 extends from the exit end of the forming assembly 110 to the outlet opening in the housing end wall 105. In FIG. 3, the guide chute 129 can be seen to be substantially rectangular in cross-section. The upstream bottom and/or side edges of the chute preferably flare outwardly to form a funnel or converging mouth inlet 160 (FIG. 5B). The top and bottom walls of the guide chute 1129 each include an opening 161 through which the wheels 132-135 extend into the interior of the guide chute (FIGS. 5A-5C). It will be appreciated that the cross-sectional dimensions (i.e., width and height) of the guide chute 129 approximate the cross-sectional dimensions of the cushioning product. The strip formed in the forming assembly 110 is urged into the guide chute 129 through its funnel inlet 160 whereat it is engaged and fed forwardly (or downstream) by the feed wheel 132 and its support wheel 133. The staggered arrangement of the flat faces 147 on the end portions 146 of the wheel 133 will cause the strip to be fed alternately from each side of its longitudinal axis, instead of just being pulled only axially. That is, the strip will be fed alternately from each side of its longitudinal axis, instead of being pulled only axially. This advance by successive pulls from one side and then the other side back and forth makes it possible to have at the center a surplus of paper with respect to its flat configuration, this surplus being generated by the rib 159 fitting in the mating groove in the wheel 132. The strip is then engaged by the compression wheel 134 and its support wheel 135. Because the wheels 134 and 135 are rotating at a slower speed than the wheels 132 and 133, the strip is longitudinally crumpled between the upstream and downstream pairs of wheels with the latter compressing folds in the strip. (For further information regarding an assembly similar to the feeding/connecting assembly 111, reference may be had to European Patent Application No. 94440027.4, filed Apr. 22, 1994 and published on Nov.2, 1995 under Publication No. 0 679 504 A1, which is hereby incorporated herein by reference.) The strip then exits the guide chute 129 and passes through the dunnage outlet opening in the end wall 105. As the strip exits the feeding/connecting assembly 111 and passes through the dunnage outlet opening in the end wall 105, the severing assembly 112 severs its leading portion into a desired length. The illustrated severing assembly 112 includes cutting components 162 preferably powered by an electric motor 163 (FIG. 1). The cutting components 162 are mounted on the downstream surface of the end wall 105 are contained within the enclosure closed by the downstream cover 107. The severing motor 163 is mounted on the base wall 103 on the side of the forming assembly opposite the feed motor 153. (See FIGS. 1 and 2.) A suitable severing assembly is disclosed in U.S. patent application Ser. No. 08/188,305, which is hereby incorporated by reference. The cut sections of dunnage then travel through the post-severing assembly 113. As seen in FIGS. 1 and 2, the post-severing assembly 113 is mounted to the downstream cover 107. The inlet and outlet of the assembly 113 are aligned with the dunnage outlet opening in the end wall 105. The post-severing assembly 113 is rectangular in cross-sectional shape and flares outwardly in the downstream direction. As the cut section of the dunnage strip, or pad, emerges from the outlet of the assembly 113, the pad is ready for use as a cushioning product. Referring now to FIGS. 17 and 18, a modified form 109 u of stock supply assembly is shown. The stock supply assembly 109, operates to layer the stock material prior to its entry into the forming assembly 110. While the stock supply assembly 109 u could be used with multi-ply stock material to double the number of layers of material, it is preferably used with single-ply stock material, in that it eliminates the need for rewinding single-ply stock material into multi-ply rolls. The stock supply assembly 109 u includes a pair of support brackets 114 u which are vertically spaced (as opposed to laterally spaced like the brackets 114) and support the stock roll R u , in a vertical orientation (the stock roll will usually be twice as wide as the normal width because the stock material is folded over on itself to provide a two layer web). The stock supply assembly 109 u further includes a layering plate 1001 which is vertically positioned upstream of the fold-down rollers 120 u , via a bracket suspending it from a pedestal on the base wall 103. The layering plate 1001 is generally triangular except that it includes a rounded entry edge 1002. As the stock material is unwound from the roll R u in a vertical plane and pulled over the layering plate 1001 into the forming assembly 110, it is folded in half into a web having two layers. This web is positioned in a horizontal plane ready for receipt by the forming assembly 110. If desired, the stock roll may be supported in a horizontal orientation with its axis oriented perpendicular to the entry path into the forming assembly 110 and an angled turner bar employed between the stock roll and the layering plate to guide the sheet material from a horizontal plane as it is payed off the stock roll to a vertical plane for passage to the layering plate 1001. It will also be appreciated that a horizontal disposition of the stock roll may also be obtained by rotating the entire machine embodiment of FIGS. 17 and 18 by 90 degrees about its longitudinal axis. In addition, additional layers may be provided by supplying stock material from one or more additional rollers, as schematically illustrated by the stock roll R v . Two, three or more stock rolls may be used with the other embodiments herein described if desired. According to another aspect of the invention, a modified version of the feeding/connecting assembly 111 may include interchangeable quick change gear sets are provided to provide respective different feed rate ratios between the input and output wheel of the wheel network. These gear sets would be similar to the gears 157-159 (FIG. 5B), except they would be of different sizes or tooth number to produce a corresponding change in feed rate ratio and thus the pad characteristics as may be desired. By employing appropriate marking on the gear sets corresponding to desired packaging applications, changes in the speed ratio could be accomplished with minimal training on the part of a machine operator by substituting the proper gear set for a given application. As explained herein, the speed ratio between the feed wheel 132 (FIG. 5C) and compression wheel 134 affects the characteristics (such as density, compactness, cushioning ability, etc.) of the pad produced during the conversion process. While the set speed ratio provided by the gear train 157-159 may be appropriate in many situations, it may be desirable to selectively change this speed ratio to alter pad characteristics Specifically, if the speed differential is increased, a stiffer, more dense pad will be produced for use in, for example, the packaging of heavier objects. On the other hand, if the speed differential is reduced, a less dense pad will be produced (possibly resulting in greater yield from a given amount of stock material) for use in, for example, the packaging of lighter objects. In another modified form of the feeding/connecting assembly, two separate feed motors could be used, one for the feed wheel shaft 137 (FIGS. 5A and 5C) and one for the compression wheel shaft 139. Either or both of the motors could have a variable speed option to allow selective adjustment of the speed ratio. It is noted that if these motors are directly coupled to the shafts 137 and 139, the need for the motion-transmitting elements 154-159 (FIG. 5A) would be eliminated. In any event, this modification would eliminate the need for the gear train 157-159 (FIG. 5A). In another modified version of the feeding/connecting assembly, shown partially in FIG. 7, the gear train 157-159 (FIG. 5A) of the drive system 128 u is replaced with a variable pitch pulley assembly 1010. In the drive system 128 u , the variable pitch pulley assembly 1010 controls the speed ratio between the feed wheel shaft 137 and the compression wheel shaft 139. The illustrated pulley 1010 includes a SL-sheave 1011 coupled to the feed wheel shaft 137, a MC-sheave 1012 coupled to the compression wheel shaft 139, and a V-belt 1013 trained therebetween. An adjustment device 1014 allows manual control (via a control knob 1015 preferably positioned outside the machine's housing for easy access) of the position of the V-belt 1013 on the sheaves 1011 and 1012 to thereby vary the speed ratio between shafts 137 and 139, in well known manner. Another modified form of the feeding/connecting assembly is shown in FIGS. 8 and 9 which is designed to provide for a convenient, and even dynamic, selective change in the biasing force between the compression wheel 134 and its support wheel 135. The support structure 129 t of the wheel network 127 t includes a pair of horizontal cross bars 131a t and 131b t which extend between, and are secured to, the side plates 130. The cross bar 131a t is vertically aligned with the shaft 138 and the cross bar 131b t is vertically aligned with the shaft 140. A first pair of pins 143a t (similar to the suspension pins 143) couple the shaft connectors 142 to the first support cross bar 131a t . The pins 143a t extend from the ends of the shaft-connectors 142 adjacent the shaft 138. Another pin 143b t is coupled to the shaft connectors 142 via a yoke 1020 connected to the ends of the shaft connectors 142 adjacent the shaft 140. The pin 143b t is attached to the cross bar 131b t via an adjustment device 1021. The adjustment device includes an adjustable stop 1021a into which the pin 143b t is threaded such that rotation of the pin will move the adjustable stop towards and away from the shaft 140. A spring 1021b t is interposed between the adjustable stop 1021a and the cross member 131b t of the yoke 1020. Accordingly, rotation of the pin will increase or decrease the biasing force acting on the yoke and in turn on the shaft 140 and wheel 135, it being noted that the pin is free to rotate relative to the yoke. As is preferred, the end of the pin projecting above the cross bar has secured thereto a knob 1022. As will be appreciated, the knob provides for easy manual adjustment of the biasing force acting on the shaft 140. The knob preferably is located external to the machine's housing, or at least at a conveniently accessible location within the machine's housing. If the knob 1022 is tightened, the biasing force between the compression wheel 134 and its support wheel 135 will be increased, thereby creating a more dense pad. If the knob 1022 is loosened, the biasing force will be decreased, thereby creating a less dense pad. Dynamic changes could be made while the machine is operating to change pad characteristics "on the fly." If desired, the knob may be replaced by other drive mechanisms, such as an electric motor that may be remotely controlled for adjustment of the biasing force. The drive system 128 w of another modified form of the feeding/connecting assembly is shown in FIG. 10. The drive system 128 w includes a reversing device 1030 which allows the reverse movement of the feeding/connecting assembly to, for example, clear paper jams in the machine. The device 1030 includes a clutch 1031 and a hand crank 1032. The clutch 1031 allows selective disengagement of the shaft of the motor 153 w from the compression wheel shaft 139. The hand crank 1032 is coupled to the compression wheel shaft 139 so that, upon disengagement of the motor drive shaft, the shaft 139 may be manually turned in the reverse direction. The hand crank 1032 can be permanently fixed to the machine as shown, or can be "folded away," or even removed during normal operation. Alternatively, the motor could be reversed to effect reverse movement of the feeding/connecting assembly. Another modified form of the feeding/connecting assembly is shown in FIGS. 20 and 21, this assembly incorporating a modified drive system 128 x . In the modified drive system 128 x , the feed wheel shaft 137 (and thus the feed wheel 132 and its support wheel 133) is directly driven by the motor 153 at a constant speed. However, the compression wheel shaft 139 (and thus the compression wheel 134 and its support wheel 135) are driven intermittently, rather than continuously, by an indexing device 1040 which replaces the gear train 157-159. When the indexed wheels 134 and 135 are not rotating, the stock material is crumpled as the rotating wheels 132 and 133 continue to advance stock material downstream. When the indexed wheels 134 and 135 are rotating, the stock material will be emitted from the feeding/connecting assembly. The indexing device 1040 is a conventional "Geneva" gear mechanism and, in the illustrated device, the compression wheel 134 rotates a quarter of a revolution for every half revolution of the feed wheel 132. The device 1040 includes a driver disk 1042 mounted to the support wall 130, a cam pin 1041 mounted to the driver disk 1042, a gear 1043 coupled to the end of the feed shaft 137, and a four-slotted disk 1044 coupled to the end of the compression wheel shaft 138. The driver disk is indexed with the compression shaft 139 so that upon every half revolution of the feed wheel shaft 137, the driver disk 1042 will also make one revolution. As the driver disk 1042 makes one revolution, it will cause the four-slotted disk 1044 to rotate a quarter of a revolution via the cam pin 1041. Another modified form 111 y of the feeding/connecting assembly is shown in FIGS. 19A-19C. The wheel network 127 y of this assembly includes a "stretching assembly" comprised of a stretch wheel 1050, its support wheel 1051, and corresponding shafts 1052 and 1053. During operation of the feeding/connecting assembly 111 y , the wheels 1050 and 1051 are rotated at a faster feed rate speed than the wheels 134 and 135 whereby the strip will be "stretched" prior to passing through the outlet opening in the end wall 105. The wheels 1050 and 1051 may be essentially identical in design and size as the wheels 134 and 135, respectively. The addition of the wheels 1050 and 1051 necessitates changes in the support structure 126 y , the wheel network 127 y , and the drive system 128 y . The support structure 126 y includes extended side walls 130 y each with an additional slot to accommodate the shaft 1053, and a cross bars 131 y positioned between each adjacent set of support wheels. In the wheel network 127 y , shaft-connectors 142 y connect all three shafts 138, 140, and 1053, and two sets of suspension pins 143 y couple the shaft-connectors 142 y to the cross bars 132 y . In the drive system 128 y , gears 1054 and 1055 are added to the gear train, gear 1054 being mounted to the stretch wheel shaft 1052 and gear 1055 being mounted to the side wall 130 y to convey motion from the gear 157 to the gear 1054. The gears 1054 and 1055 may be sized so that the stretch wheel 1050 is rotated anywhere between a feed rate speed just slightly faster than the compression wheel 134 to a feed rate speed equal to the feed wheel 132. Also, although not shown in FIGS. 19A-19C, the guide chute 129 (FIGS. 5A-5C) is preferably elongated and its slots modified to accommodate the wheels 1050 and 1051. In a further modified form 111 z of the feeding/connecting assembly shown in FIGS. 22-24, a movable barrier 1060 replaces the compression wheel 134, its support wheel 135, and the compression wheel shaft 139. The barrier 1060 is spring biased towards the feed wheel 132 so that as the strip of cushioning is expelled therefrom, it will be restricted by the barrier 1060, thereby crumpling the strip in a longitudinal direction. As pressure applied by the crumpling strip increases, the spring bias of the barrier 1060 will be overcome, and it will open to allow the crumpled strip to pass through the outlet opening in the end wall 105. The illustrated barrier 1060 is made from a circular (in cross-section) bar formed into a rectangular loop having rounded corners. The loop is perpendicularly bent at a central portion to form a rounded corner 1061 between an upper portion 1062 and a lower portion 1063 of the barrier 1060. The corner 1061 of the barrier 1060 is rotatably attached around the shaft 140 (previously used for the support wheel 135). When in a rest position, the barrier's lower portion 1063 extends into the guide chute 129 z in a downward and downstream sloping direction with its upper portion 1062 extending upwardly therefrom. In the wheel network 127 z , a guide pin 1064 is connected to, and extends horizontally from, cross bar 131. The pin 1064 is attached at its other end to a bracket 1065 secured to the top portion 1062 of the barrier, and a spring 1064a is carried on the pin 1064 and interposed between the bracket 1065 and the cross bar 131. As the pressure of the crumpling strip increases behind the lower portion 1063 of the barrier, the upper portion of the barrier 1062 will be pushed towards the cross-bar 131 thereby pivoting the lower portion 1063 upward to allow release of the strip. In the guide chute 129 z , the upper slot 161 z is extended to the downstream edge of the guide chute, which extends beyond the outlet opening in the end wall 105. (See FIG. 22.) The drive system 128 z is essentially the same as the drive system 128, except that the gear train 157-159 is eliminated. In FIGS. 6A and 6B, a cushioning conversion machine 200 is shown. The machine 200 converts sheet-like stock material into a three-dimensional cushioning product of a desired length. As with the machine 100, the preferred stock material for the machine 200 consists of plural plies or layers of biodegradable and recyclable sheet-like stock material such as 30 to 50 pound Kraft paper rolled onto a hollow cylindrical tube to form a roll R of the stock material. However, the stock material would preferably consist of three plies of paper and, in any event, would not be intermittently glued together. As with the machine 100, the preferred cushioning product of the machine 200 has lateral accordion-like or pillow-like portions and is connected, or assembled, along a relatively thin central band separating the pillow-like portions. The machine 200 is similar to the machine 100 discussed above, and includes an essentially identical housing 202, feeding/connecting assembly 211, severing assembly 212, and post-severing assembly 213. However, the stock supply assembly 209 and the forming assembly 210 of the machine 200 differ from these assemblies in the machine 100. The stock supply assembly 209 includes two support brackets 214 which are laterally spaced apart and mounted to the machine's frame, or more particularly the upstream wall (or rectangular border) 208. The stock supply assembly 209 also includes a sheet separator 216, and a constant-entry roller 218. The sheet separator 216 includes three vertically spaced rollers which extend between, and are connected to, the support brackets 214. (The number of separator rollers corresponds to the number of plies or layers of the stock material whereby more or less rollers could be used depending on the number of layers.) The constant-entry roller 218 also extends between, and is connected to, the support brackets 214. As the paper is unwound from the supply roll R, it travels over the constant-entry roller 218 and into the separating device 216. In the separating device, the plies or layers of the stock material are separated by the separator rollers and this "pre-separation" is believed to improve the resiliency of the produced cushioning product. The constant-entry roller 218 provides a non-varying point of entry for the stock material into the separator 216 regardless of the diameter of the roll R. (Details of a similar stock supply assembly are set forth in U.S. Pat. No. 5,322,477, the entire disclosure of which is hereby incorporated by reference.) The forming assembly 210 includes a shaping chute 219 and a forming member 220. The shaping chute 219 is longitudinally converging in the downstream direction and is positioned in a downstream portion of the enclosure formed by the machine's housing. Its entrance is outwardly flared in a trumpet-like fashion and its exit is positioned adjacent the feeding/connecting assembly 211. The chute 219 is mounted to the housing at the bottom wall 103 and at 221. The forming member 220 has a "pinched U" or "bobby pin" shape including a bight portion joining upper and lower legs. The lower leg extends to a point approximately coterminous with the exit end of the shaping chute 219. The rearward portion of the forming member 220 preferably projects rearwardly of the entry end of the shaping chute by approximately one-half its overall length. Also, the radius of the rounded base or bight portion is approximately one-half the height of the mouth of the shaping chute. This provides for a smooth transition from the separating device 216 to the forming member and then into the shaping chute. The lower leg 220a of the forming member 220 extends generally parallel to the bottom wall 219a of the shaping chute 219. However, the relative inclination and spacing between the lower leg of the forming member and bottom wall of the shaping chute may be adjusted as needed to obtain proper shaping and forming of the lateral edges of the stock material. Such adjustment may be effected and then maintained by an adjustment device 223 which, as best shown in FIG. 6C, extends between the legs of the forming member at a point midway along the length of the lower leg, it being noted that the upper leg may be shorter as only sufficient length is needed to provide for attachment of the top wall of the shaping chute. The adjustment device 223 includes a rod 224 having a lower end attached to the lower leg of the forming member 220 by a rotation joint 225 (such as a ball-and-socket joint). The upper threaded end of the rod 224 extends through a threaded hole in the top wall of the shaping chute as well as through a threaded hole in a upper leg of the forming member 220 and is held in place by a nut 224a secured to the shaping chute 219. To adjust the gap between the lower leg of the forming member and the bottom wall of the shaping chute, the top of the threaded rod is turned the appropriate direction. The rod's top may be provided with a screwdriver slot or wrench flats, to easily accomplish this turning with standard tools. Further details of the preferred chute 219 and shaping member 220 are set forth in U.S. application Ser. No. 08/487,182, the entire disclosure of which is hereby incorporated by reference. However, it should be noted that other chutes and shaping members are possible with, and contemplated by, the present invention. By way of example, the chutes and/or shaping members set forth in U.S. Pat. Nos. 4,026,198; 4,085,662; 4,109,040; 4,717,613; and 4,750,896, could be substituted for the forming chute 219 and/or the shaping member 220. As the stock material passes through the shaping chute 219, its lateral end sections are rolled or folded inwardly into generally spiral form and are urged inwardly toward one another so that the inwardly rolled edges form a pillow-like portions of stock material disposed in lateral abutting relationship as they emerge from the exit end of the shaping chute. The forming member 220 coacts with the shaping chute 219 to ensure proper shaping and forming of the paper, the forming ember being operative to guide the central section of the stock material along the bottom wall of the chute 219 for controlled inward rolling of the lateral side sections of the stock material. The rolled stock material, or strip, then travels to the feeding/connecting assembly 211. Another cushioning conversion machine 300, formed from modular units 300a and 300b according to the present invention, is shown in FIGS. 11A, 11B,11C and 12. The machine 300 converts sheet-like stock material into a three-dimensional cushioning product of a desired length. As with the machines 100 and 200, the preferred cushioning product of the machine 300 has lateral crumpled pillow-like portions and is connected, or assembled, along a central band separating the pillow-like portions. As with the machines 100 and 200, the preferred stock material for the machine 300 consists of plural plies or layers of biodegradable and recyclable sheet-like stock material such as 30 to 50 pound Kraft paper rolled onto a hollow cylindrical tube to form a roll R of the stock material. The first modular unit 300a includes a housing 302a similar to the downstream portion of the housing 102 of the machine 100. (See FIG. 11A.) A feeding/connecting assembly 311, a severing assembly 312 and a post-severing assembly 313, which are essentially identical to the corresponding assemblies in the machine 100, are mounted to the housing 302a in the same manner as they are mounted the downstream portion of the housing 102. However, an expanding device 370 occupies the space in the machine housing 102 that had been occupied by the forming assembly 110 and requires less space. (See FIG. 11A.) Additionally, a guide roller 372 is mounted to the upstream end of the housing 302a via brackets 374. The expanding device 370 includes a mounting member 378 to which a separating member 380 is joined. (See FIGS. 11B and 11C.) The mounting member 378 includes a transverse support or mounting arm 381 having an outwardly turned end portion 383 and an oppositely turned end portion 385 to which the separating member 380 is attached. The outer end portion 383 is mounted to the housing 302a by a bracket 387 and suitable fastening elements. The separating member 380 includes a transverse support 393 and fold expansion elements 395 at opposite ends of the transverse support 393 that are relatively thicker than the transverse support 393, with respect to the narrow dimension of the stock material. In the illustrated expanding device, the mounting member 378 is formed by a rod or tube, and the fold expansion elements are formed by rollers supported for rotation on the transverse support at opposite ends thereof. The transverse support 393 is attached near one end thereof to the adjacent end portion 385 of mounting member 381 for support in cantilevered fashion. The expanding device 373 is designed for use with flat-folded stock material which is formed by the second modular unit 300b. During the conversion process, the layers of the stock material (formed by the edge and central portions of the ply or plies) travel through the expanding device 373. More particularly, the central section of the folded stock material travels over the sides of the rollers 395 opposite the mounting arm 381, while the inner edge portion of the stock material travels in the narrow V-shape or U-shape slot formed between the transverse support 393 and the mounting arm 381 and the other or outer edge portion of the travels over the side of the mounting arm 381 furthest the separating member 380. As a result, the lateral end sections are separated from one another and from the central section, thereby introducing loft into the then expanded material which now takes on a three dimensional shape as it enters the guide chute of the feeding/connecting device 311. Further details of the expanding device 370 are set forth in U.S. patent application Ser. No. 08/584,092, which is hereby incorporated herein by reference in its entirety. The second modular unit 300b includes a housing 302b similar to the upstream portion of the housing 102 of the machine 100. (See FIG. 12.) A forming assembly 310 is essentially identical to, and is mounted to the housing 302b in the same manner as, the corresponding assembly in the machine 100. However, a stock roll R may be supported by a floor mounted stand or stock roll support 2002. Additionally, a guide roller 398 is mounted to a downstream end of the housing 302a via bracket 399. A packaging system 2000 incorporating the cushioning conversion machine 300 is shown in FIG. 13. In addition to the machine 300, the system includes a table 2001 and a floor-mounted stock support 2002. The first modular unit 300a is located on top of the table 2001 and the second modular unit 300b is located below the table. As the stock material is unwound from the roll R, it travels from the support 2002, over the plate 119 through the forming assembly 310, under the guide roller 398 (positioned between the legs of the table), over the guide roller 372, through the expanding device 370 and into the feeding/connecting assembly 311. The strip is then severed by the severing assembly 312 and the cut section travels through the post-severing assembly 313. A modified version 2000 u of the packaging system is shown in FIG. 14. In the packaging system 2000 u , the folded stock material from the unit 300b passes through an opening 2003 in the table 2001 u . This arrangement allows a more central positioning of the units 300a and 300b relative to the table 2001 u and also protects the folded strip from interference as it travels between the units. Another modified version 2000 w of the packaging system is shown in FIG. 15. In the packaging system 2000 w , the first unit 300a is stacked on top of the second unit 300b below an elevated (when compared to tables 2001 and 2001 w ) table 2001 w . Additionally, the post-severing assembly 313 w is curved upwardly towards an opening 2003 w in the table whereby the cut section of cushioning will be deposited on the table top. This arrangement allows the table top to be clear of all machine components during the production of cushioning products. Another packaging system 2000 x according to the present invention is shown in FIG. 16. This packaging system incorporates a machine 300 x which is similar to the machine 300 except for its first modular unit 300a x . Specifically, the unit 300a x has manual, rather than motor-powered, severing assembly 312 x . Additionally, the housing 300b x is in the form of a two part casing. The other components, such as the expanding device 370 and the feeding/connecting assembly 311, operate in essentially the same manner as described above. For further details of the unit 300b x , reference may be had to U.S. patent application Ser. No. 08/584,092. One may now appreciate that the present invention provides an improved cushioning conversion machine related methodology. Although the invention has been shown and described with respect to certain preferred embodiments, it is obvious that equivalent alterations and modifications will occur to others skilled in the art upon the reading and understanding of this specification. The present invention includes all such equivalent alterations and modifications. Accordingly, while a particular feature of the invention may have been described above with respect to only one of the illustrated embodiments, such feature may be combined with one or more features of the other embodiments, as may be desired and advantageous for any given or particular application. It is noted that the position references in the specification (i.e, top, bottom, lower, upper, etc.) are used only for ease in explanation when describing the illustrated embodiments and are in no way intended to limit the present invention to particular orientation. Also, the terms (including a reference to a "means") used to identify the herein-described assemblies and devices are intended to correspond, unless otherwise indicated, to any assembly/device which performs the specified function of such an assembly/device that is functionally equivalent even though not structurally equivalent to the disclosed structure which performs the function in the illustrated exemplary embodiment of the invention.

An improved cushioning conversion machine and related methodology characterized by one or more features including, inter alia, a feeding/connecting assembly which enables an operator to easily vary a characteristic, for example, the density, of the cushioning product; a feeding/connecting assembly wherein input and/or output wheels or rollers thereof are made at least in part of an elastomeric or other friction enhancing material, which reduces the cost and complexity of the input and output rollers; a manual reversing mechanism that is useful, for example, for clearing paper jams; a modular arrangement of a forming assembly and feeding/connecting assembly in separate units that may be positioned remotely from one another, as may be desired for more efficient utilization of floor space; a turner bar which enables alternative positioning of a stock supply roll; and a volume expanding arrangement cooperative with the feeding/connecting assembly for adjusting the density of the cushioning product and changing product yield. The features of the invention may be individually or collectively used in cushioning conversion machines of various types.

BACKGROUND OF THE INVENTION This invention relates to a carrier suitable for lifting a block of material, and more particularly to a carrier for lifting a block of salt. One particularly arduous task is the handling of large salt blocks. It is always desirable to simplify work that must be done. These blocks are handled by cattle feeders and by home owners. The cattle feeders handle the blocks because the salt is a necessary, nutritional element in the raising of the cattle. It is always desirable to put out the salt blocks for the cattle to lick. To do this requires lifting and handling of salt blocks. The homeowner uses the salt blocks in a water softener. The salt block also is the more efficient way to use salt in a water softener. Nevertheless, the salt block is not favored by the homeowner because of its bulkiness and its difficulty to handle. It thus becomes desirable to provide a device permitting simplified handling of salt blocks. Grippers are known for grabbing blocks of material. Some of these grippers are too sharp and cut the salt block when used therewith. Other grippers do not have a simplified release mechanism. Still other grippers are heavy and cumbersome to use. Such grippers do not overcome the problems or simplify and reduce the effort required to lift or otherwise maneuver a salt block. The desirability of developing a device suitable for lifting or otherwise handling a salt block thus becomes extremely clear. If a lighter weight gripper can be successfully developed, sufficient strength must be maintained to handle the weight of the salt block. It is difficult to maintain the strength to handle the weight of the salt block while at the same time providing a device which can be easily released from the salt block. Thus, problems abound in developing a device to handle a salt block. SUMMARY OF THE INVENTION Therefore, it is an object of this invention to provide a device for lifting a salt block. A further object of this invention is to provide a device which minimizes damage to the salt block. A still further object of this invention is to provide a device which is easily released from a salt block. Yet a further object of this invention is to provide a device which is more simple to operate. Also an object of this invention is to provide a device which is not cumbersome to operate. Another object of this invention is to provide a light weight device capable of lifting a salt block. These and other objects of the invention are met by providing a device for lifting a salt block including a handle and pivotally mounted arms designed to fit the salt block with blunt gripping knobs. BRIEF DESCRIPTION OF THE DRAWING FIG. I is perspective view of the front of salt block carrier 10 in contact with salt block 12. FIG. II is a plan, reverse view of FIG. I. FIG. III is an end view of FIG. II. FIG. IV is a top view of FIG. III. FIG. V is a bottom view of FIG. II. Throughout the Figures of drawing where the same part appears in more than one Figure of the drawing, the same numeral is applied thereto. DESCRIPTION OF THE PREFERRED EMBODIMENTS A salt block carrier has a handle having two arms secured thereto at one end and having gripping knobs at the other end of each arm, which knobs fit onto a salt block; and is easily released from the salt block as desired. Referring now to FIGS. I and II, salt block carrier 10 is depicted in contact with a salt block 12. Standardly, salt blocks such as salt block 10 are symmetrical and have salt indentations 14. The salt indentation 14 is matched by an equivalent indentation on the side of the salt block 12 oppositely disposed therefrom and not visible. Each of the indentations 14 is a recess into the block 12 and is a standard feature of salt block 12 commonly used in water softeners and cattle feeding. Indentations 14 provide for an easy mold release after the salt block has been molded or otherwise shaped. Indentations 14 receive salt block carrier 10 and provide for gripping of the salt block 12. Salt block carrier 10 includes handle 20, and first arm 60 and second arm 100 secured to handle 20. Handle 20 has a grip area 22 which is generally the portion of handle 20 in contact with the hand of a person using the salt block carrier 10. Adjacent grip area 22 at one end thereof is a first handle side 24. Opposite first handle side 24 and substantially symmetrical thereto is a second handle side 26. Grip area 22 naturally extends into first handle side 24 and second handle side 26 as a matter of course. In a preferred fashion, first handle side 24 and second handle side 26 are generally in a arcuate right angle position relative to grip area 22. First handle base 28 is an extension of first handle side 24 and joins second handle base 30. Second handle base 30 is clearly an extension of second handle side 26. Thus, in the preferred form as shown handle 20 is of a generally pentagonal shape with two adjoining arcuate right angles. Handle interior 32 is generally the area through which the hand is inserted for the purpose of operating the salt block carrier 10. Handle exterior 34 is, of course, the oppositely disposed and outside portion of handle 20. In a preferred design form, handle 20 includes an interior flange 36 on handle interior 32 in a flattened and rounded edge shape. Likewise, exterior flange 38 is similar to interior flange 36 and oppositely disposed therefrom. Thus, exterior flange 38 is at handle exterior 34. Flange base 40 is generally between interior flange 36 and exterior flange 38 in a substantially perpendicular fashion such that an I-beam type shape is achieved. This structure provides strength for handle 20, and provides for comfort in lifting because of the flattened and rounded-edge nature of handle interior 32. On handle exterior 34 and first handle base 28 is a first base arm mount 42. Similarly on handle exterior 34 and second handle base 30 is a second base arm mount 44. First base arm mount 42 and second base arm mount 44 are substantially symmetrical. Both first base arm mount 42 and second arm mount 44 have a mount base 46. Each mount base 46 is secured by molding or other suitable fashion to the respective first handle base 28 and second handle base 30. Mount top 48 is oppositely disposed from mount base 46 and gives first base arm mount 42 and second base arm mount 44 its generally quadrilateral shape. Within first base arm mount 42 is a first base aperture (not shown). Similarly, a second base aperture (not shown) is within second base arm mount 44. First arm rod 54 passes through the first base aperture to secure first arm 60 thereto. By the same token second arm rod 56 passes through the second base aperture to secure second arm 100 to handle 20. First arm 60 includes a first apertured end 62 having a first arm aperture 64. First arm aperture 64 protrudes through first arm base 66 and first arm holder 68 to thereby form first arm slot 70. First arm holder 68 and first arm base 66 provide first arm slot 70 for receiving first base arm mount 42. First arm aperture 64 aligns with the first base aperture between first arm holder 68 and first arm base 66. First arm rod 54 is then slidably and frictionally inserted through first arm aperture 64 and the first base aperture to secure first arm 60 to handle 20 at first base arm mount 42. First arm 60 further includes adjacent to first arm slot 66 a first arm groove 72. First arm groove 72 extends throughout first arm 60 down to first arm grip 80. First arm valley 74 is similar to first arm groove 72 down the opposite side of first arm 60. First arm holder 68 is on the same side as first arm groove 72. First arm base 66 is on the same side as first arm valley 74. However, first arm valley 74 is a mirror image of first arm groove 72. First arm arch 76 permits first apertured end 62 to extend into first arm side 78. Generally first arm arch 76 is a 90 degree arch permitting first arm side 78 to be perpendicular to first apertured end 62. First arm side 78 extends downwardly into first arm grip 80. First arm grip 80 is at an obtuse angle with first arm side 78. First arm grip 80 includes a first grip base 82 and a first grip knob 84. First grip base 82 is the extension at the obtuse angle of first arm side 78. First grip knob 84 is a rounded and enlarged end of first arm grip 80. The rounded and enlarged first grip knob 84 permits gripping of the salt block 12 in the salt indentations 14. On first apertured end 62 is a first arm release 86. First arm release 86 is triangular in shape and has a first arm vertex 88 extending out from first apertured end 62. First arm 60 includes a first arm interior 90 within which the salt block is held. Oppositely disposed from first arm interior 90 is first arm exterior 92. First arm exterior 92 includes first arm exterior flange 94. First arm flange 94 permits a flattened appearance for first arm 60 similar to the flange of handle 20 caused by exterior flange 38. First arm interior 90 includes first arm interior flange 96 which is flattened therein on the interior 90 of salt block carrier 10. First flange connector 98 connects first arm exterior flange 94 and first arm interior flange 96 to form an I beam style cross-section. In this fashion, a decorative, strong and useful salt block carrier 10 is depicted. The first arm release 86 contacts the salt block 12 when the handle 20 has approached the salt block 12 due to downward pressure on the handle 20. As the handle 20 is pressed downwardly, first arm release 86 contacts the salt block forcing first arm 60 to rotate about first arm rod 54 in an upward direction thereby releasing salt block carrier 10 from the salt block 12. Second arm 100 functions similarly. Second arm 100 includes a second apertured end 102 having a second arm aperture 104. Second arm aperture 104 protrudes through second arm base 106 and second arm holder 108 to thereby form second arm slot 110. Second arm holder 108 and second arm base 106 provide the second arm slot 110 for receiving second base arm mount 44. Second arm aperture 104 aligns with the second base aperture in second base arm mount 44 between second arm holder 108 and second arm base 106. Second arm rod 56 is then slidably and frictionally inserted through second arm aperture 104 and the second base aperture to secure second arm 100 to handle 20 at second base arm mount 44. Second arm 100 further includes adjacent to second arm slot 110 a second arm groove 112. Second arm groove 112 extends throughout second arm 100 down to second arm grip 120. Second arm valley 114 is similar to second arm groove 112 down the opposite side of second arm 100. Second arm holder 108 is on the same side as second arm groove 112. Second arm base 106 is on the same side as second arm valley 114. However, second arm valley 114 is a mirror image of second arm groove 112. Second arm arch 116 permits second apertured end 102 to extend into second arm side 118. Generally second arm arch 116 is a 90 degree arch permitting second arm side 118 to be perpendicular to second apertured end 102. Second arm side 118 extends downwardly into second arm grip 120. Second arm grip 120 is at an obtuse angle with second arm side 118. Second arm grip 120 includes a second grip base 122 and a second grip knob 124. Second grip base 122 is the extension at the obtuse angle of second arm side 118. Second grip knob 124 is a rounded and enlarged end of second arm grip 120. The rounded and enlarged knob permits gripping of the salt block 12 in the salt indentations 14. On second apertured end 102 is a second arm release 126. Second arm release 126 is triangular in shape and has a second arm vertex 128 extending out from second apertured end 102. Second arm 100 includes a second arm interior 130 within which the salt block is held. Oppositely disposed from second arm interior 130 is second arm exterior 132. Second arm exterior 132 includes second arm exterior flange 134. Second arm exterior flange 134 permits a flattened appearance for second arm 100 similar to the flange of handle 20 caused by exterior flange 38. Second arm interior 130 includes second arm interior flange 136 which is flattened therein on the interior 130 of salt block carrier 10. Second flange connector 138 connects second arm exterior flange 134 and second arm interior flange 136 to form an I beamed style cross-section. In this fashion, a decorative, strong and useful salt block carrier 10 is depicted. The second arm release 126 contacts the salt block 12 when the handle has approached the salt block 12 due to downward pressure on the handle 20. As the handle 20 is pressed downwardly, second arm release 126 contacts the salt block 12 forcing second arm 100 to rotate about second arm rod 56 in an upward direction thereby releasing salt block carrier 10 from the salt block 12. Referring now to FIG. III, handle 20 is shown from a side view. First base arm mount 42 is depicted as a narrower portion of handle 20. First arm 60 is connected thereto. Referring now to FIG. IV, a top view of salt bock carrier 10 is depicted. In this fashion, handle 20 and more particularly grip 22 of handle 20 is shown along with first arm 60 and second arm 100. More particularly first apertured end 62 of first arm 60 is shown in a top view and second apertured end 102 is shown in a top view. It is clear that first arm 60 and second arm 100 are symmetrical. It is also clear that salt block carrier 10 has an axis of symmetry along Line 2--2. Referring now to FIG. V, a bottom view of salt block carrier 10 is depicted. From the bottom view, handle 20 is seen and more particularly first handle base 28 and second handle base 30. Also, a bottom view of the salt block first arm 60 and second arm 100 are seen. Clearly visible also are the first arm release 86 and the second arm release 126. In this fashion, the structure of the salt block carrier is depicted. In operation, salt block carrier 20 is placed to show that first grip knob 84 and second grip knob 124 contact oppositely disposed salt block indentations 14. First arm 60 and second arm 100 are pressed so that first grip knob 84 and second grip knob 124 are within the indentations. Due to the presence of first grip knob 84 and second grip knob 124, lifting of the salt block 12 by handle 20 becomes feasible. The first grip knob 84 and second grip knob 124 cause the gripping of salt block 12 as handle 20 is pulled then in an upward direction at grip area 22. In order to release salt block carrier 20, one merely needs to press down on handle 20 in order to have first arm release 86 and second arm release 126 contact the salt block and push first arm 60 and second arm 100 in an outward direction. In this manner, salt block carrier 20 is released from the salt block 10. Any material having the desired shaping characteristics and lightweight strength capabilities may be used to form salt block carrier 10 of this invention. A preferred material is a moldable plastic or synthetic resin having the required strength and stiffness characteristic for the desired purpose. Typically, a moldable thermosettable resin is used. Even a metal or metallic alloy may be used if the desired lightweight strength characteristic can be met. Combinations of metal and resin or plastic may also be used.

A salt block carrier includes a handle with a pair of pivotally mounted arms movably secured to the handle. The arms each have an end designed to fit and hold the salt block with blunt gripping knobs.

BACKGROUND OF THE INVENTION This invention relates to a four wheel drive system for a vehicle, wherein power circulation is prevented, thereby eliminating excessive loads, worn tires and deteriorated fuel consumption and the like due to power circulation. Recently, as the drive systems used in vehicles such as motor vehicles, four wheel drive systems have become popular. Heretofore, there have been two types of such four wheel drive systems. The first is one having a two wheel - four wheel switching control clutch, wherein front wheels and rear wheels are selectively driven and connected to each other. The second is one having a center differential gear and a differential control clutch for restricting (up to the degree of prohibiting) the differential action of the center differential gear, whereby the front wheels and the rear wheels can be directly or almost directly connected to each other. These are disclosed in bulletins of Japanese Patent Kokai (Laid-Open) Nos. 72420/1980, 53520/1983 and 101829/1983. When the above-described release and restriction of differential action between front wheels and rear wheels are controlled in response to an outside signal, an arrangement may be adopted in which the differential action is restricted throughout the running time of the vehicle or during most of the running time, and the restriction of differential action is suitably released (or the degree of the restriction changed) in accordance with the running conditions of the vehicle. Or, an arrangement may be adopted in which the differential action between the front wheels and the rear wheels is permitted during normal operation, and the restriction of differential action is applied when necessary. In general, when the difference in rotary speed between front wheels and rear wheels is increased, it is recognized that either the front wheels or the rear wheels are slipping relative to the road surface, whereupon the differential action is restricted. However, conventional four wheel drive systems have not taken into account differences in effective radii of tires caused by changed in weight and load acting on front wheels or rear wheels, change in pneumatic pressures within the tires, changing wear in the tires, use of a temper tire, attachment of chains to the tires, and the like. Accordingly, if the restriction of differential action is performed in the case when the such differences exist in the effective radii of the tires, power circulation occurs, whereby excessive load, abrasion of tires, deterioration in fuel consumption and the like may be induced. SUMMARY OF THE INVENTION The present invention has been developed to obviate the above-described disadvantages of the prior art and has as its object the provision of a four wheel drive system for a vehicle, such that when there is a difference in effective radii of the tires, and the restriction of differential action would cause excessive load and abrasion of tires due to the power circulation, the above-described disadvantages are prevented from occurring. To achieve the above-described object, as shown in FIG. 1 (A), a first aspect of the present invention contemplates a four wheel drive system for a vehicle, wherein the restriction of the differential action between front wheels and rear wheels is released or the degree of restriction is changed as necessary. The system includes: means for detecting whether or not the running conditions of the vehicle are stable; means for temporarily releasing the restriction of differential action between front wheels and rear wheels when the running conditions of the vehicle are detected to be stable; means for detecting a differential rate between front wheels and rear wheels upon release of the restriction of differential action; and means for prohibiting (inclusive of concept of restricting) the restriction of differential action when the differential rate is detected to be eaual to a predetermined value or thereabove. Furthermore, to achieve the above-described object, as shown in FIG. 1 (B), a second aspect of the present invention is directed to a four wheel drive system for a vehicle, wherein the differential action between front wheels and rear wheels is restricted as necessary. The system includes: means for detecting whether or not the running conditions of the vehicle are stable; means for detecting whether or not differential action between front wheels and rear wheels is being restricted; means for detecting a differential rate between front wheels and rear wheels when it is detected that the running conditions of the vehicle are stable and the differential action between front wheels and the rear wheels is not being restricted; and means for prohibiting (inclusive of cencept of restricting) the restriction of differential action when the differential rate is greater than or equal to a predetermined value. The first aspect of the present invention relates to a four wheel drive system for a vehicle, constructed such that the differential action between front wheels and rear wheels is restricted throughout the running time of the vehicle or during most of the running time, and the restriction is released or the degree of the restriction is changed as necessary. There are four types of four wheel drive systems wherein the differential action between front wheels and rear wheels is released or restricted. These include: (1) those wherein the drive can be switched between two wheel and four wheel by an ON-OFF control clutch, (2) those wherein the drive can be switched between two wheel and four wheel by a control clutch having a variable transmitting capacity such as a wet-type multi-disk clutch, (3) those wherein a center differential gear, positioned between the front wheels and rear wheels, is prohibited or released by an ON-OFF control clutch, and (4) those wherein a center differential gear, positioned between the front wheels and rear wheels, is controlled by a control clutch having a variable transmitting capcity such as a wet-type multi-disk clutch: and, in this invention, any type of the above may be adopted. In order to detect the differential rate between front wheels and rear wheels, differential action is permitted. According to a first aspect of the present invention, since the differential action between front wheels and rear wheels is restricted throughout the running time of the vehicle or during most of the running time, in order to detect the differential rate, the restriction of differential action is temporarily released. One of the features of the first aspect of the present invention resides in this active release of the restriction of differential action for the purpose of detecting the differential rate. The temporary release of the restriction of differential action is performed only when the running conditions of the vehicle are determined to be stable. As a result, the differential rate between front wheels and rear wheels can be detected under conditions where the change in behavior of the vehicle caused by release of restriction of differential action is extremely small, and a small differential rate caused by a difference in the effective radii of the tires can be easily and accurately detected. When the differential rate is detected to be greater than or equal to the predetermined value by this detection, the restriction of differential action thenceforth is prohibited (or restricted), whereby front wheels and rear wheels can rotate in accordance with the difference in effective radii of the tires, so that the excessive loads, unusual wears, deterioration in fuel consumption and the like due to power circulation can be prevented. The second aspect of the present invention relates to a four wheel drive system for a vehicle, constructed such that the differential action between front wheels and rear wheels is permitted during normal operation, and the differential action therebetween is restricted as necessary. The second aspect of the present invention is similar to the first aspect in that the specific type or types of four wheel drive systems are not limited. In the second aspect of the present invention, the differential rate between front wheels and rear wheels detected when the running conditions are detected to be stable and the restriction of differential action is not performed. By so doing, a small differential rate caused by a difference in effective radius of the tires can be easily and accurately detected, and a change in the behavior of the vehicle due to detection of the differential rate can be completely prevented. More specifically, in the four wheel drive system according to the second aspect of the present invention, differential action between front wheels and rear wheels is permitted during most of the running time of the vehicle, whereby this time may be utilized to detected the differential rate. When the differential rate is detected to be greater than or equal to a predetermined value, the restriction of differential action thenceforth is prohibited (or restricted), so that excessive loads, unusual wears of the tires, deterioration in fuel consumption and the like due to power circulation can be prevented similarly to the first aspect of the present invention. According to the first and second aspects of the present invention, a preferably embodiment is an arrangement by which an apparatus for detecting whether or not the running conditions of the vehicle are stable detects whether or not such conditions are established, i.e., whether the vehicle speed is greater than or equal to a predetermined value, the steering angle is less than or equal to a predetermined value and the engine load is less than or equal to a predetermined value. As a result, the differential rate can be detected while the vehicle is in a stable operating condition. Furthermore, according to the first aspect of the present invention, a preferable embodiment is an arrangement such that, when the differential rate is detected to be smaller than a predetermined value, the release of the restriction of differential action for the detection of the differential rate is stopped until the vehicle stops. With this arrangement, frequent detection of the differential rate is discontinued, so that the original control of the restriction of differential action in accordance with the various running conditions is fully performed. Furthermore, according to the first aspect of the present invention, a preferable embodiment is an arrangement such that, when the differential rate is detected to be greater than or equal to a predetermined value, restriction of differential action is prohibited, and detection of the differential rate is repeated each time the running conditions of the vehicle are detected to be stable, until the vehicle stopped. With this arrangement, prohibition of restriction of differential action due to a mistaken detection can be prevented. More specifically, even after restriction of the differential is prohibited, if the differential rate is detected to be smaller than the predetermined value a few times, for example, the prohibition of restriction of differential action may be released. When starting from a stop, before the conditions for detecting the differential rate are established, the values at the time of the previous detection may be stored and used. With this arrangement, when the value of the differential rate of the previous time is smaller than a predetermined value, the original restriction is maintained until the value is detected again. When the value of the differential rate of the previous time is greater than or equal to a predetermined value, differential restriction can be prohibited until the value is detected again. BRIEF DESCRIPTION OF THE DRAWINGS The above object, features and advantages of the present invention, as well as other objects and advantages thereof, will become more apparent from the description of the invention which follows, taken in conjunction with the accompanying drawings, wherein like reference characters designate the same to similar parts and wherein: FIGS. 1 (A) and 1 (B) are block diagrams showing technical illustration of the first and second aspects of the present invention, respectively; FIG. 2 is a conceptual diagram showing the power transmission system of the four wheel drive system for the vehicle, to which the present invention is applied; FIG. 3 is a conceptual diagram of hydraulic control circuit of the differential control clutch of the system shown in FIG. 2; FIG. 4 is a flow chart showing the control procedure used in the four wheel drive system of FIG. 2; and FIG. 5 is a flow chart showing another control procedure. DESCRIPTION OF THE PREFERRED EMBODIMENT The present invention will be described in detail with reference to the accompanying drawings which illustrate preferred embodiments of the present invention. Referring to FIG. 2, the four wheel drive systems includes an engine 10, anautomatic transmission 20, a center differential gear unit 30, a front differential gear unit 40, a transfer unit 50, a rear differential gear unit 60, a differential control clutch unit 70, a control unit 80, and various input systems 90. The engine 10 is laterally placed in the front portion of the vehicle. Output from the engine is transmitted to the automatic transmission 20. The automatic transmission 20 includes a fluid type torque converter 21 andan auxiliary transmission 22. A hydraulic control device 23 shifts the automatic transmission 20 among four forward gears and one backward gear. The hydraulic control device 23 is controlled in response to a original from the control unit 80. The power through the automatic transmission 20 is transmitted through an output gear 24 to an input gear 31 which is partof the center differential gear unit 30. The center differential gear unit 30 includes a differential case 32 for integrally supporting the input gear 31, two differential pinions 34 and 35 positioned opposite one another and rotatably supported by a pinion shaft 33 secured to the differential case 32, a side gear 36 for transmitting power to the rear wheels, and a side gear 37 for transmittingpower to the front wheels. The side gear 36 is meshed with the differentialpinions 34 and 35, and is connected to a transfer ring gear 51 which is part of transfer unit 50. The side gear 37 is meshed with the differentialpinions 34 and 35, and is connected through a hollow front wheel drive shaft 41 to a differential case 42 which is part of the front differentialgear unit 40. The front differential gear unit 40 includes two differential pinions 44 and 45 positioned opposite one another and rotatably supported by a pinionshaft 43 secured to the differential case 42, a side gear 46 for transmitting power to the left front wheel, and a side gear 47 for transmitting power to the right front wheel. The side gear 46 and 47 are each meshed with the differential pinions 44 and 45. A shaft 48 of the left front wheel is connected to the side gear 46 for transmitting power to the left front wheel. And a shaft 49 of the right front wheel is connected to the side gear 47 for transmitting power to the right front wheel. The transfer unit 50 includes a transfer ring gear 51 connected to the sidegear 36 for transmitting power to the rear wheels, a driven pinion 52 meshing with the transfer ring gear 51, and an output rotary gear 54 whichis integrally rotatable with the driven pinion 52 through a propeller shaft53. The output rotary gear 54 is connected to the rear differential gear unit 60. The rear differential gear unit 60 includes a differential case 61 formed with a ring gear which is meshed with the output rotary gear 54, two differential pinions 63 and 64 which are positioned opposite one another and which are rotatably supported by a pinion shaft 62 secured to the differential case 61, a side gear 65 for transmitting power to the left rear wheel, and a side gear 66 for transmitting power to the right rear wheel. The side gear 65 is meshed with the differential pinions 63 and 64,and is connected to a shaft 67 if the left rear wheel, and the side gear 66is meshed with the differential pinions 63 and 64, and is connected to a shaft 68 of the right rear wheel. The differential control clutch unit 70 is adapted to selectively connect the differential case 32, which is an input member of the center differential gear 30, to the front wheel drive shaft 41, which is an output member of the center differential gear 30. The clutch unit 70 includes a wet-type multi-disk clutch 71 and a hydraulic control device 72for controlling the multi-disk clutch 71. More specifically, as shown in FIG. 3, the multi-disk clutch 71 is provided with a hydraulic servo device73. A servo piston 75 is urged to the right against the resilient force of a return spring 76 by servo oil pressure fed to an oil chamber 74 of the hydraulic servo device 73. With this arrangement, the differential case 32and the front wheel drive shaft 41 are connected to each other, and the torque transmitting capacity is increased or decreased in proportion to the increase or decrease of the servo oil pressure fed to the oil chamber 74. The hydraulic control device 72 supplies the servo oil pressure to theoil chamber 74 of the hydraulic servo device 73. The hydraulic control device 72 includes a line oil pressure control valve 77 for regulating oilpressure of an oil pump 81 incorporated in the automatic transmission 20. The pressure of the oil supplied by the oil pump corresponds to the engineload, and is controlled by an electromagnetic servo hydraulic control valve78. The servo hydraulic control valve 78 includes a port "a" connected to the oil chamber 74, an oil pressure port "b", to which the line oil pressure from the line oil pressure control valve 77 is supplied, and a drain port "c". The servo hydraulic control valve 78 connects the port "a"to the oil pressure port "b" while current passed through the control valve78 and connects the port "a" to the drain port "c" while current is not passed through the control valve 78. The servo hydraulic control valve 78 is controlled in response to pulse signals from the control device 80. Thepulse signals have a specific duty ratio. With this arrangement, servo oil pressure commensurate to the duty ratio is fed to the oil chamber 74. The control unit 80 controls the hydraulic control devices 23 and 72 in response to respective input signals from the input systems 90. The control unit 80 receives throttle opening information from a throttle opening sensor 91, manual shift range (position of manual shift lever) information of the automatic transmission 20 from a manual shift position sensor 92, front wheel rotary speed information from front wheel rotary speed sensors 93 and 94, rear wheel rotary speed information from a rear wheel rotary speed sensor 95, vehicle steering angle information from a steering angle sensor 96. Other sensors may be included described. Upon receiving these input signals, the control unit 80 outputs control signalsfor controlling gear stages of the automatic transmission 20 to the hydraulic control device 23 in accordance with the manual range information, the front wheel rotary speed or the rear wheel rotary speed (vehicle speed), and the throttle opening. Furthermore, the control device80 outputs pulse signals having the specific duty raito for controlling thetorque transmitting capacity of the differential control clutch 70 to the hydraulic control device 72 in accordance with input torque to the automatic transmission 20 and the gear stages. FIG. 4 shows the inerative control procedure in the above system. The interative control procedure is an embodiment of the first aspect of the present invention. In order to detect the differential rate Rd between thefront wheels and the rear wheels, it is detected whether the vehicle speed v is greater than or equal to a predetermined value V1, whether the steering angle Ast is less than or equal to predetermined value Ast1, and whether the throttle opening Ath is less than or equal to a predetermined value Ath1. Such conditions prevail when carrying a light load running steady in a straight line. When the above conditions are detected, the restriction of differential action is released. Accodringly, the differential rate Rd can be detected since the difference between the behavior of the vehicle when the prevailing restriction of differential action and when operating free of differential action is extremely small, and since a very small differential rate due to the difference in effective radius of the tires can be easily and accurately detected. Furthermore, in the control procedure, when the differential rate Rd is detected to be smaller than a predetermined value Rd1, the differential rate Rd is not detected again before the vehicle is stopped. When this occurs, frequent detections of the differential rate Rd (which requires frequent release of restriction of differential action), is discontinued, so that original control of the restriction of differential action in accordance with the running conditions can be fully performed. On the contrary, so long as the differential rate Rd is detected to be greater than or equal to a predetermined value, the differential rate Rd is detected every time the detecting conditions are established. In this case, if the differential rate is detected to be smaller than a predetermined value N consecutive times, restriction of differential action is permitted again. This is because the detection of the differential rate may have been imprecise, and original restriction of differential action should be performed repeatedly. Until the detecting conditions are established for the first time after restart, the result of the detection of the previous time is stored and used as current value. Accordingly, when the result of the previous time is smaller than the predetermined value, the original restriction of differential action is performed under the conditions of the restriction of differential action corresponding to the various running conditions until the detecting conditions are again established. When the previous results are greater than or equal to the predetermined value, the restriction of differential action is prohibited until the detecting conditions are again established. The flow chart shown in FIG. 4 will hereunder be described in detail. In Step 201, whether or not a flag F is set is detected. The flag F is adapted to indicate whether or not the differential rate Rd has been detected since the last time the vehicle was brought to a stop. When F is equal to 1, it is indicated that the detection has been performed. When F is not equal to 1, it is indicated that the detection has not been performed since the last stop. When F is equal to 1, the routine proceeds to Step 209, and, when F is not equal to 1, the routine proceeds to Step 202. In Step 202, whether or not vehicle speed v is greater than or equal to a predetermined value V1, i.e., whether v≧V1 is detected. When v≧V1, the routine proceeds to Step 203. When v<V1, the routine proceeds to Step 208. In Step 208, whether or not the differential rate Rdis less than or equal to a predetermined value Rd1, i.e., whether Rd≦Rd1 is detected. In this case, i.e., when the answer to Step 202is "no", the differential rate Rd is the value of the previous detecting time, which has been stored. In Step 203, whether or not steering angle Ast is less than or equal to thepredetermined value Ast1, i.e., whether Ast≦Ast1 is detected. When Ast≦Ast1, the routine proceeds to Step 204. When Ast>Ast1, the routine proceeds to Step 208 such that a new value for Rd is not calculated. In Step 204, whether or not throttle opening Ath is less than or equal to the predetermined throttle opening Ath1, i.e., whether Ath≦Ath1, isdetected. When Ath≦Ath1, the routine proceeds to Step 205, when Ath>Ath1, the routine proceeds to Step 208 such that a new value for Rd isnot calculated. When the routine proceeds to Step 205, the current detecting conditions areestablished, i.e., the vehicle running conditions are determined to be stable, and the flag F is set at 1. In Step 206, the torque transmitting capacity Tc of the restriction of the differential control clutch unit 70 is set at 0, and the center differential gear unit 30 is made free i.e., the center differential is unrestricted. In Step 207, the rotary numbers of the left and right front wheels and the left and right rear wheels are detected, whereby the mean rotary speed of the front wheels nF and the mean rotary speed of the rear wheels nR are calculated, and the differential rate Rd is calculated in accordance with the following equation. Rd=|1-(nF / nR)| (1) In Step 208, whether or not the differential rate Rd is smaller than the predetermined value Rd1 is detected. When Rd<Rd1, the original control of the restriction of differential action is performed thereafter. When Rd≧Rd1, the routine proceeds to Step 210, where the restriction of differential action is prohibited regardless of the running conditions thenceforth. In Step 211, a counter C is set at a predetermined value N (N=3 for example). In Step 209, whether or not the differential rate Rd, detected the previoustime is smaller than the predetermined value Rd1 is detected. When Rd<Rd1, the process goes out of this routine and the original control of the restriction of differential action is performed. When Rd≧Rd1, the routine proceeds to Step 212 in order to detect the differential rate Rd again. In Step 212, whether or not the vehicle speed v≧V1 is detected. When v≧V1, the routine proceeds to Step 213. When v<V1, the process goes out of this routine and prohibition of the restriction ofdifferential action is continued. The value of the counter C is not changed. In Step 213, whether or not the steering angle Ast≦Ast1 isdetected. When Ast≦Ast1, the routine proceeds to Step 214. When Ast>Ast1, the process goes out of this routine. In Step 214, whether or not the throttle opening Ath≦Ath1 is detected. When Ath≦Ath1, the routine proceeds to Step 215. When Ath>Ath1, the process goes out of this routine. In Step 215, the differential rate Rd iscalculated according to equation (1) above. In Step 216, whether or not thedifferential rate Rd<Rd1 is detected. When Rd<Rd1, the routine proceeds to Step 217. When Rd≧Rd1, the routine proceeds to Step 210. Namely, when Rd≧Rd1, the routine goes through Steps 210 and 211, and the counter C is set at N again. In Step 217, the counter value C is reduced by 1. In Step 218, whether or not C is equal to 0 is determined. When C isnot equal to 0, the process goes out of this routine and, when C is equal to 0, the routine proceeds to Step 219. In Step 219, because it has been determined N consecutive times that the differential rate Rd is smaller than the predetermined value Rd1, the restriction of differential action is permitted, whereby the original control of the restriction of differential action is performed. The second embodiment of the present invention will hereunder be described in connection with in FIG. 5. The above-described first embodiment is of the type wherein differential action of the center differential gear is constantly restricted (the four wheel drive system according to the first aspect of the present invention). However, in a second embodiment in accordance with the presentinvention, such an arrangement is adopted using a four wheel drive system similar to that of the first embodiment, wherein the center differential gear is held free during the normal time, and the differential action is suitably restricted in accordance with the running conditions of the vehicle, for example, during slippage of the front or rear wheels or when accelerating from a stop. FIG. 5 shows an example of the control procedure adopted in the above-described second embodiment of the four wheel drive system. In FIG. 5, Steps of the control procedure similar to those shown in FIG. 4 are designated by like reference numerals. As is apparent from FIG. 5, similarly to the preceding embodiment, in Steps 202-204, whether or not the running conditions are stable is detected, and thereafter, in Step 300, whether or not the differential action of the center differential gear 30 is currently free or not is detected. When the running conditions are stable and the differential action is free, the routine proceeds to Step 207, where the differential rate Rd is calculated. In Step 208, whether or not the differential rate Rd thus calculated is smaller than the predetermined value Rd1 is detected. When Rd<Rd1, the routine proceedsto Step 302, where the original control of differential action is performed. When Rd≧Rd1, the routine proceeds to Step 304 and the restriction of differential action is prohibited thenceforth. In the four wheel drive system wherein the center differential gear is madefree during normal operation as described above, because the original restriction of differential action can be performed without hindrance, thedifferential rate Rd can be calculated whenever the conditions necessary for calculating the differential rate Rd are established, whereby the differential rate Rd can be assigned an updated value.

In a four wheel drive vehicle, wherein either a state where differential action between front wheels and rear wheels is restricted or a permitted state can be selected, the differential rate is calculated during stable running conditions, and, when the differential rate is greater than or equal to a predetermined value, the restriction of differential action is prohibited. As a result, power circulation can be prevented when there is a difference in effective radius between the front wheels and rear wheels due to, e.g., use of a temper tire, mounting of chains on the tires, or the like.

BACKGROUND [0001] 1. Field [0002] Certain aspects of the present disclosure generally relate to neural system engineering and, more particularly, to developing and testing of a neural network by improving the execution time of a parameter search. [0003] 2. Background [0004] Artificial neural networks may provide innovative and useful computational techniques for certain applications in which traditional computational techniques are cumbersome, impractical, or inadequate. Artificial neural networks may have corresponding structure and/or function in biological neural networks. An artificial neural network, which may be an interconnected group of artificial neurons (i.e., neuron models), is a computational device or represents a method to be performed by a computational device. Because artificial neural networks can infer a function from observations, such networks are particularly useful in applications where the complexity of the task or data makes the design of the function by conventional techniques burdensome. SUMMARY [0005] In one aspect, a method of wireless communication is disclosed. The method includes serializing sub-systems of the system by determining the one-way dependencies between the sub-systems and/or parallelizing the sub-systems by determining independencies within each sub-system. The method further includes pruning input parameters of each sub-system based on whether each input parameter affects each sub-system. [0006] Another aspect discloses an apparatus including means for serializing sub-systems of the system by determining the one-way dependencies between the sub-systems and/or means for parallelizing the sub-systems by determining independencies within each sub-system. The apparatus further includes means for pruning input parameters of each sub-system based on whether each input parameter affects each sub-system. [0007] In another aspect, a computer program product for wireless communications in a wireless network having a non-transitory computer-readable medium is disclosed. The computer readable medium has non-transitory program code recorded thereon which, when executed by the processor(s), causes the processor(s) to perform operations of serializing sub-systems of the system by determining the one-way dependencies between the sub-systems and/or parallelizing the sub-systems by determining independencies within each sub-system. The program code also causes the processor(s) to prune input parameters of each sub-system based on whether each input parameter affects each sub-system. [0008] Another aspect discloses wireless communication having a memory and at least one processor coupled to the memory. The processor(s) is configured to serialize sub-systems of the system by determining the one-way dependencies between the sub-systems and/or parallelize the sub-systems by determining independencies within each sub-system. The processor(s) is also configured to prune input parameters of each sub-system based on whether each input parameter affects each sub-system. [0009] Additional features and advantages of the disclosure will be described below. It should be appreciated by those skilled in the art that this disclosure may be readily utilized as a basis for modifying or designing other structures for carrying out the same purposes of the present disclosure. It should also be realized by those skilled in the art that such equivalent constructions do not depart from the teachings of the disclosure as set forth in the appended claims. The novel features, which are believed to be characteristic of the disclosure, both as to its organization and method of operation, together with further objects and advantages, will be better understood from the following description when considered in connection with the accompanying figures. It is to be expressly understood, however, that each of the figures is provided for the purpose of illustration and description only and is not intended as a definition of the limits of the present disclosure. BRIEF DESCRIPTION OF THE DRAWINGS [0010] The features, nature, and advantages of the present disclosure will become more apparent from the detailed description set forth below when taken in conjunction with the drawings in which like reference characters identify correspondingly throughout. [0011] FIG. 1 illustrates an example network of neurons in accordance with certain aspects of the present disclosure. [0012] FIG. 2 illustrates an example of a processing unit (neuron) of a computational network (neural system or neural network) in accordance with certain aspects of the present disclosure. [0013] FIG. 3 illustrates an example of spike-timing dependent plasticity (STDP) curve in accordance with certain aspects of the present disclosure. [0014] FIG. 4 illustrates an example of a positive regime and a negative regime for defining behavior of a neuron model in accordance with certain aspects of the present disclosure. [0015] FIGS. 5A-5B illustrate examples of a neural network according to aspects of the present disclosure. [0016] FIGS. 6A-6D illustrate examples of a neural network according to aspects of the present disclosure. [0017] FIG. 7 is a block diagram illustrating a method for improving a parameter evaluation of a neural network according to an aspect of the present disclosure. [0018] FIG. 8 is a block diagram illustrating functions for improving a parameter evaluation of a neural network according to an aspect of the present disclosure. [0019] FIG. 9 illustrates an example implementation of designing a neural network using a general-purpose processor in accordance with certain aspects of the present disclosure. [0020] FIG. 10 illustrates an example implementation of designing a neural network where a memory may be interfaced with individual distributed processing units in accordance with certain aspects of the present disclosure. [0021] FIG. 11 illustrates an example implementation of designing a neural network based on distributed memories and distributed processing units in accordance with certain aspects of the present disclosure. [0022] FIG. 12 illustrates an example implementation of a neural network in accordance with certain aspects of the present disclosure. [0023] FIG. 13 is a block diagram illustrating a method for performing a parameter sweep over a system having sub-systems with one-way dependencies DETAILED DESCRIPTION [0024] The detailed description set forth below, in connection with the appended drawings, is intended as a description of various configurations and is not intended to represent the only configurations in which the concepts described herein may be practiced. The detailed description includes specific details for the purpose of providing a thorough understanding of the various concepts. However, it will be apparent to those skilled in the art that these concepts may be practiced without these specific details. In some instances, well-known structures and components are shown in block diagram form in order to avoid obscuring such concepts. [0025] Based on the teachings, one skilled in the art should appreciate that the scope of the disclosure is intended to cover any aspect of the disclosure, whether implemented independently of or combined with any other aspect of the disclosure. For example, an apparatus may be implemented or a method may be practiced using any number of the aspects set forth. In addition, the scope of the disclosure is intended to cover such an apparatus or method practiced using other structure, functionality, or structure and functionality in addition to or other than the various aspects of the disclosure set forth. It should be understood that any aspect of the disclosure disclosed may be embodied by one or more elements of a claim. [0026] The word “exemplary” is used herein to mean “serving as an example, instance, or illustration.” Any aspect described herein as “exemplary” is not necessarily to be construed as preferred or advantageous over other aspects. [0027] Although particular aspects are described herein, many variations and permutations of these aspects fall within the scope of the disclosure. Although some benefits and advantages of the preferred aspects are mentioned, the scope of the disclosure is not intended to be limited to particular benefits, uses or objectives. Rather, aspects of the disclosure are intended to be broadly applicable to different technologies, system configurations, networks and protocols, some of which are illustrated by way of example in the figures and in the following description of the preferred aspects. The detailed description and drawings are merely illustrative of the disclosure rather than limiting, the scope of the disclosure being defined by the appended claims and equivalents thereof. An Example Neural System, Training and Operation [0028] FIG. 1 illustrates an example artificial neural system 100 with multiple levels of neurons in accordance with certain aspects of the present disclosure. The neural system 100 may have a level of neurons 102 connected to another level of neurons 106 through a network of synaptic connections 104 (i.e., feed-forward connections). For simplicity, only two levels of neurons are illustrated in FIG. 1 , although fewer or more levels of neurons may exist in a neural system. It should be noted that some of the neurons may connect to other neurons of the same layer through lateral connections. Furthermore, some of the neurons may connect back to a neuron of a previous layer through feedback connections. [0029] As illustrated in FIG. 1 , each neuron in the level 102 may receive an input signal 108 that may be generated by neurons of a previous level (not shown in FIG. 1 ). The signal 108 may represent an input current of the level 102 neuron. This current may be accumulated on the neuron membrane to charge a membrane potential. When the membrane potential reaches its threshold value, the neuron may fire and generate an output spike to be transferred to the next level of neurons (e.g., the level 106 ). In some cases, the neuron may continuously transmit a signal to the next level of neurons. The transmitted signal may be a function of the membrane potential. Such behavior can be emulated or simulated in hardware and/or software, including analog and digital implementations such as those described below. [0030] In biological neurons, the output spike generated when a neuron fires is referred to as an action potential. This electrical signal is a relatively rapid, transient, nerve impulse, having an amplitude of roughly 100 mV and a duration of about 1 ms. In a particular embodiment of a neural system having a series of connected neurons (e.g., the transfer of spikes from one level of neurons to another in FIG. 1 ), every action potential has basically the same amplitude and duration, and thus, the information in the signal may be represented only by the frequency and number of spikes, or the time of spikes, rather than by the amplitude. The information carried by an action potential may be determined by the spike, the neuron that spiked, and the time of the spike relative to other spike or spikes. The importance of the spike may be determined by a weight applied to a connection between neurons, as explained below. [0031] The transfer of spikes from one level of neurons to another may be achieved through the network of synaptic connections (or simply “synapses”) 104 , as illustrated in FIG. 1 . Relative to the synapses 104 , neurons of level 102 may be considered pre-synaptic neurons and neurons of level 106 may be considered post-synaptic neurons. The synapses 104 may receive output signals (i.e., spikes) from the level 102 neurons and scale those signals according to adjustable synaptic weights w 1 (i,i+1) , . . . , w P (i,i+1) where P is a total number of synaptic connections between the neurons of levels 102 and 106 and i is an indicator of the neuron level. For example, in the example of FIG. 1 , i represents neuron level 102 and i+1 represents neuron level 106 . Further, the scaled signals may be combined as an input signal of each neuron in the level 106 . Every neuron in the level 106 may generate output spikes 110 based on the corresponding combined input signal. The output spikes 110 may be transferred to another level of neurons using another network of synaptic connections (not shown in FIG. 1 ). [0032] Biological synapses can mediate either excitatory or inhibitory (hyperpolarizing) actions in postsynaptic neurons and can also serve to amplify neuronal signals. Excitatory signals depolarize the membrane potential (i.e., increase the membrane potential with respect to the resting potential). If enough excitatory signals are received within a certain time period to depolarize the membrane potential above a threshold, an action potential occurs in the postsynaptic neuron. In contrast, inhibitory signals generally hyperpolarize (i.e., lower) the membrane potential. Inhibitory signals, if strong enough, can counteract the sum of excitatory signals and prevent the membrane potential from reaching a threshold. In addition to counteracting synaptic excitation, synaptic inhibition can exert powerful control over spontaneously active neurons. A spontaneously active neuron refers to a neuron that spikes without further input, for example due to its dynamics or a feedback. By suppressing the spontaneous generation of action potentials in these neurons, synaptic inhibition can shape the pattern of firing in a neuron, which is generally referred to as sculpturing. The various synapses 104 may act as any combination of excitatory or inhibitory synapses, depending on the behavior desired. [0033] The neural system 100 may be emulated by a general purpose processor, a digital signal processor (DSP), an application specific integrated circuit (ASIC), a field programmable gate array (FPGA) or other programmable logic device (PLD), discrete gate or transistor logic, discrete hardware components, a software module executed by a processor, or any combination thereof. The neural system 100 may be utilized in a large range of applications, such as image and pattern recognition, machine learning, motor control, and alike. Each neuron in the neural system 100 may be implemented as a neuron circuit. The neuron membrane charged to the threshold value initiating the output spike may be implemented, for example, as a capacitor that integrates an electrical current flowing through it. [0034] In an aspect, the capacitor may be eliminated as the electrical current integrating device of the neuron circuit, and a smaller memristor element may be used in its place. This approach may be applied in neuron circuits, as well as in various other applications where bulky capacitors are utilized as electrical current integrators. In addition, each of the synapses 104 may be implemented based on a memristor element, where synaptic weight changes may relate to changes of the memristor resistance. With nanometer feature-sized memristors, the area of a neuron circuit and synapses may be substantially reduced, which may make implementation of a large-scale neural system hardware implementation more practical. [0035] Functionality of a neural processor that emulates the neural system 100 may depend on weights of synaptic connections, which may control strengths of connections between neurons. The synaptic weights may be stored in a non-volatile memory in order to preserve functionality of the processor after being powered down. In an aspect, the synaptic weight memory may be implemented on a separate external chip from the main neural processor chip. The synaptic weight memory may be packaged separately from the neural processor chip as a replaceable memory card. This may provide diverse functionalities to the neural processor, where a particular functionality may be based on synaptic weights stored in a memory card currently attached to the neural processor. [0036] FIG. 2 illustrates an exemplary diagram 200 of a processing unit (e.g., a neuron or neuron circuit) 202 of a computational network (e.g., a neural system or a neural network) in accordance with certain aspects of the present disclosure. For example, the neuron 202 may correspond to any of the neurons of levels 102 and 106 from FIG. 1 . The neuron 202 may receive multiple input signals 204 1 - 204 N (X 1 9 N ), which may be signals external to the neural system, or signals generated by other neurons of the same neural system, or both. The input signal may be a current, a conductance, a voltage, a real-valued, and/or a complex-valued. The input signal may comprise a numerical value with a fixed-point or a floating-point representation. These input signals may be delivered to the neuron 202 through synaptic connections that scale the signals according to adjustable synaptic weights 206 1 - 206 N (W 1 -W N ), where N may be a total number of input connections of the neuron 202 . [0037] The neuron 202 may combine the scaled input signals and use the combined scaled inputs to generate an output signal 208 (i.e., a signal Y). The output signal 208 may be a current, a conductance, a voltage, a real-valued, and/or a complex-valued. The output signal may be a numerical value with a fixed-point or a floating-point representation. The output signal 208 may be then transferred as an input signal to other neurons of the same neural system, or as an input signal to the same neuron 202 , or as an output of the neural system. [0038] The processing unit (neuron) 202 may be emulated by an electrical circuit, and its input and output connections may be emulated by electrical connections with synaptic circuits. The processing unit 202 and its input and output connections may also be emulated by a software code. The processing unit 202 may also be emulated by an electric circuit, whereas its input and output connections may be emulated by a software code. In an aspect, the processing unit 202 in the computational network may be an analog electrical circuit. In another aspect, the processing unit 202 may be a digital electrical circuit. In yet another aspect, the processing unit 202 may be a mixed-signal electrical circuit with both analog and digital components. The computational network may include processing units in any of the aforementioned forms. The computational network (neural system or neural network) using such processing units may be utilized in a large range of applications, such as image and pattern recognition, machine learning, motor control, and the like. [0039] During the course of training a neural network, synaptic weights (e.g., the weights w 1 (i,i+1) , . . . , w P (i,i+1) from FIG. 1 and/or the weights 206 1 - 206 N from FIG. 2 ) may be initialized with random values and increased or decreased according to a learning rule. Those skilled in the art will appreciate that examples of the learning rule include, but are not limited to the spike-timing-dependent plasticity (STDP) learning rule, the Hebb rule, the Oja rule, the Bienenstock-Copper-Munro (BCM) rule, etc. In certain aspects, the weights may settle or converge to one of two values (i.e., a bimodal distribution of weights). This effect can be utilized to reduce the number of bits for each synaptic weight, increase the speed of reading and writing from/to a memory storing the synaptic weights, and to reduce power and/or processor consumption of the synaptic memory. Synapse Type [0040] In hardware and software models of neural networks, the processing of synapse related functions can be based on synaptic type. Synapse types may include non-plastic synapses (no changes of weight and delay), plastic synapses (weight may change), structural delay plastic synapses (weight and delay may change), fully plastic synapses (weight, delay and connectivity may change), and variations thereupon (e.g., delay may change, but no change in weight or connectivity). The advantage of multiple types is that processing can be subdivided. For example, non-plastic synapses may not execute plasticity functions (or wait for such functions to complete). Similarly, delay and weight plasticity may be subdivided into operations that may operate together or separately, in sequence or in parallel. Different types of synapses may have different lookup tables or formulas and parameters for each of the different plasticity types that apply. Thus, the methods would access the relevant tables, formulas, or parameters for the synapse's type. Use of varying synapse types may add flexibility and configurability to an artificial neural network. [0041] There are implications of spike-timing dependent structural plasticity being executed independently of synaptic plasticity. Structural plasticity may be executed even if there is no change to weight magnitude (e.g., if the weight has reached a minimum or maximum value, or it is not changed due to some other reason) because structural plasticity (i.e., an amount of delay change) may be a direct function of pre-post spike time difference. Alternatively, structural plasticity may be set as a function of the weight change amount or based on conditions relating to bounds of the weights or weight changes. For example, a synapse delay may change only when a weight change occurs or if weights reach zero but not if they are at a maximum value. However, it may be advantageous to have independent functions so that these processes can be parallelized reducing the number and overlap of memory accesses. Determination of Synaptic Plasticity [0042] Neuroplasticity (or simply “plasticity”) is the capacity of neurons and neural networks in the brain to change their synaptic connections and behavior in response to new information, sensory stimulation, development, damage, or dysfunction. Plasticity is important to learning and memory in biology, as well as for computational neuroscience and neural networks. Various forms of plasticity have been studied, such as synaptic plasticity (e.g., according to the Hebbian theory), spike-timing-dependent plasticity (STDP), non-synaptic plasticity, activity-dependent plasticity, structural plasticity and homeostatic plasticity. [0043] STDP is a learning process that adjusts the strength of synaptic connections between neurons. The connection strengths are adjusted based on the relative timing of a particular neuron's output and received input spikes (i.e., action potentials). Under the STDP process, long-term potentiation (LTP) may occur if an input spike to a certain neuron tends, on average, to occur immediately before that neuron's output spike. Then, that particular input is made somewhat stronger. On the other hand, long-term depression (LTD) may occur if an input spike tends, on average, to occur immediately after an output spike. Then, that particular input is made somewhat weaker, and hence the name “spike-timing-dependent plasticity.” Consequently, inputs that might be the cause of the post-synaptic neuron's excitation are made even more likely to contribute in the future, whereas inputs that are not the cause of the post-synaptic spike are made less likely to contribute in the future. The process continues until a subset of the initial set of connections remains, while the influence of all others is reduced to an insignificant level. [0044] Because a neuron generally produces an output spike when many of its inputs occur within a brief period, (i.e., inputs being sufficiently cumulative to cause the output), the subset of inputs that typically remains includes those that tended to be correlated in time. In addition, because the inputs that occur before the output spike are strengthened, the inputs that provide the earliest sufficiently cumulative indication of correlation will eventually become the final input to the neuron. [0045] The STDP learning rule may effectively adapt a synaptic weight of a synapse connecting a pre-synaptic neuron to a post-synaptic neuron as a function of time difference between spike time t pre of the pre-synaptic neuron and spike time t post of the post-synaptic neuron (i.e., t=t post −t pre ). A typical formulation of the STDP is to increase the synaptic weight (i.e., potentiate the synapse) if the time difference is positive (the pre-synaptic neuron fires before the post-synaptic neuron), and decrease the synaptic weight (i.e., depress the synapse) if the time difference is negative (the post-synaptic neuron fires before the pre-synaptic neuron). [0046] In the STDP process, a change of the synaptic weight over time may be typically achieved using an exponential decay, as given by: [0000] Δ   w  ( t ) = { a +   - t / k + + μ , t > 0 a -   t / k - , t > 0 , ( 1 ) [0000] where k + and k − τ sign(Δt) are time constants for positive and negative time difference, respectively, a + and a − are corresponding scaling magnitudes, and μ is an offset that may be applied to the positive time difference and/or the negative time difference. [0047] FIG. 3 illustrates an exemplary diagram 300 of a synaptic weight change as a function of relative timing of pre-synaptic and post-synaptic spikes in accordance with the STDP. If a pre-synaptic neuron fires before a post-synaptic neuron, then a corresponding synaptic weight may be increased, as illustrated in a portion 302 of the graph 300 . This weight increase can be referred to as an LTP of the synapse. It can be observed from the graph portion 302 that the amount of LTP may decrease roughly exponentially as a function of the difference between pre-synaptic and post-synaptic spike times. The reverse order of firing may reduce the synaptic weight, as illustrated in a portion 304 of the graph 300 , causing an LTD of the synapse. [0048] As illustrated in the graph 300 in FIG. 3 , a negative offset u may be applied to the LTP (causal) portion 302 of the STDP graph. A point of cross-over 306 of the x-axis (y=0) may be configured to coincide with the maximum time lag for considering correlation for causal inputs from layer i−1. In the case of a frame-based input (i.e., an input that is in the form of a frame of a particular duration of spikes or pulses), the offset value μ can be computed to reflect the frame boundary. A first input spike (pulse) in the frame may be considered to decay over time either as modeled by a post-synaptic potential directly or in terms of the effect on neural state. If a second input spike (pulse) in the frame is considered correlated or relevant to a particular time frame, then the relevant times before and after the frame may be separated at that time frame boundary and treated differently in plasticity terms by offsetting one or more parts of the STDP curve such that the value in the relevant times may be different (e.g., negative for greater than one frame and positive for less than one frame). For example, the negative offset μ may be set to offset LTP such that the curve actually goes below zero at a pre-post time greater than the frame time and it is thus part of LTD instead of LTP. Neuron Models and Operation [0049] There are some general principles for designing a useful spiking neuron model. A good neuron model may have rich potential behavior in terms of two computational regimes: coincidence detection and functional computation. Moreover, a good neuron model should have two elements to allow temporal coding. For example, the arrival time of inputs affects output time and coincidence detection can have a narrow time window. Additionally, to be computationally attractive, a good neuron model may have a closed-form solution in continuous time and stable behavior including near attractors and saddle points. In other words, a useful neuron model is one that is practical and that can be used to model rich, realistic and biologically-consistent behaviors, as well as be used to both engineer and reverse engineer neural circuits. [0050] A neuron model may depend on events, such as an input arrival, output spike or other event whether internal or external. To achieve a rich behavioral repertoire, a state machine that can exhibit complex behaviors may be desired. If the occurrence of an event itself, separate from the input contribution (if any), can influence the state machine and constrain dynamics subsequent to the event, then the future state of the system is not only a function of a state and input, but rather a function of a state, event, and input. [0051] In an aspect, a neuron n may be modeled as a spiking leaky-integrate-and-fire neuron with a membrane voltage v n (t) governed by the following dynamics: [0000]  v n  ( t )  t = α   v n  ( t ) + β  ∑ m  w m , n  y m  ( t - Δ   t m , n ) , ( 2 ) [0000] where α and β are parameters, w m,n w m,n is a synaptic weight for the synapse connecting a pre-synaptic neuron m to a post-synaptic neuron n, and y m (t) is the spiking output of the neuron m that may be delayed by dendritic or axonal delay according to Δt m,n until arrival at the neuron n's soma. [0052] It should be noted that there is a delay from the time when sufficient input to a post-synaptic neuron is established until the time when the post-synaptic neuron actually fires. In a dynamic spiking neuron model, such as Izhikevich's simple model, a time delay may be incurred if there is a difference between a depolarization threshold v t and a peak spike voltage v peak . For example, in the simple model, neuron soma dynamics can be governed by the pair of differential equations for voltage and recovery, i.e.: [0000]  v  t = ( k  ( v - v t )  ( v - v r ) - u + I ) / C , ( 3 )  u  t = a  ( b  ( v - v r ) - u ) . ( 4 ) [0000] where v is a membrane potential, u is a membrane recovery variable, k is a parameter that describes time scale of the membrane potential v, a is a parameter that describes time scale of the recovery variable u, b is a parameter that describes sensitivity of the recovery variable u to the sub-threshold fluctuations of the membrane potential v, v r is a membrane resting potential, I is a synaptic current, and C is a membrane's capacitance. In accordance with this model, the neuron is defined to spike when v>v peak . Hunzinger Cold Model [0053] The Hunzinger Cold neuron model is a minimal dual-regime spiking linear dynamical model that can reproduce a rich variety of neural behaviors. The model's one- or two-dimensional linear dynamics can have two regimes, wherein the time constant (and coupling) can depend on the regime. In the sub-threshold regime, the time constant, negative by convention, represents leaky channel dynamics generally acting to return a cell to rest in a biologically-consistent linear fashion. The time constant in the supra-threshold regime, positive by convention, reflects anti-leaky channel dynamics generally driving a cell to spike while incurring latency in spike-generation. [0054] As illustrated in FIG. 4 , the dynamics of the model 400 may be divided into two (or more) regimes. These regimes may be called the negative regime 402 (also interchangeably referred to as the leaky-integrate-and-fire (LIF) regime (which is different from the LIF neuron model)) and the positive regime 404 (also interchangeably referred to as the anti-leaky-integrate-and-fire (ALIF) regime, not to be confused with the ALIF neuron model)). In the negative regime 402 , the state tends toward rest (v − ) at the time of a future event. In this negative regime, the model generally exhibits temporal input detection properties and other sub-threshold behavior. In the positive regime 404 , the state tends toward a spiking event (v s ). In this positive regime, the model exhibits computational properties, such as incurring a latency to spike depending on subsequent input events. Formulation of dynamics in terms of events and separation of the dynamics into these two regimes are fundamental characteristics of the model. [0055] Linear dual-regime bi-dimensional dynamics (for states v and u) may be defined by convention as: [0000] τ p   v  t = v + q ρ ( 5 ) - τ u   u  t = u + r ( 6 ) [0000] where q p and r are the linear transformation variables for coupling. [0056] The symbol ρ is used herein to denote the dynamics regime with the convention to replace the symbol ρ with the sign “−” or “+” for the negative and positive regimes, respectively, when discussing or expressing a relation for a specific regime. [0057] The model state is defined by a membrane potential (voltage) v and recovery current u. In basic form, the regime is essentially determined by the model state. There are subtle, but important aspects of the precise and general definition, but for the moment, consider the model to be in the positive regime 404 if the voltage v is above a threshold (v + ) and otherwise in the negative regime 402 . [0058] The regime-dependent time constants include τ − which is the negative regime time constant, and τ + which is the positive regime time constant. The recovery current time constants τ u is typically independent of regime. For convenience, the negative regime time constant τ − is typically specified as a negative quantity to reflect decay so that the same expression for voltage evolution may be used as for the positive regime in which the exponent and τ + will generally be positive, as will be τ u . [0059] The dynamics of the two state elements may be coupled at events by transformations offsetting the states from their null-clines, where the transformation variables are: [0000] q p =−τ ρ βu−v ρ   (7) [0000] r =δ( v +ε)  (8) [0000] where δ, ε, β and v − , v + are parameters. The two values for v ρ are the base for reference voltages for the two regimes. The parameter v − is the base voltage for the negative regime, and the membrane potential will generally decay toward v − in the negative regime. The parameter v + is the base voltage for the positive regime, and the membrane potential will generally tend away from in the positive regime. [0060] The null-clines for v and u are given by the negative of the transformation variables q p and r, respectively. The parameter δ is a scale factor controlling the slope of the u null-cline. The parameter ε is typically set equal to −v − . The parameter β is a resistance value controlling the slope of the v null-clines in both regimes. The τ ρ time-constant parameters control not only the exponential decays, but also the null-cline slopes in each regime separately. [0061] The model may be defined to spike when the voltage v reaches a value v S . Subsequently, the state may be reset at a reset event (which may be one and the same as the spike event): [0000] v={circumflex over (v)} −   (9) [0000] u=u+Δu   (10) [0000] where {circumflex over (v)} − and Δu are parameters. The reset voltage {circumflex over (v)} − is typically set to v − . [0062] By a principle of momentary coupling, a closed form solution is possible not only for state (and with a single exponential term), but also for the time required to reach a particular state. The close form state solutions are: [0000] v  ( t + Δ   t ) = ( v  ( t ) + q ρ )   Δ   t τ ρ - q ρ ( 11 ) u  ( t + Δ   t ) = ( u  ( t ) + r )   - Δ   t τ u - r ( 12 ) [0063] Therefore, the model state may be updated only upon events, such as an input (pre-synaptic spike) or output (post-synaptic spike). Operations may also be performed at any particular time (whether or not there is input or output). [0064] Moreover, by the momentary coupling principle, the time of a post-synaptic spike may be anticipated so the time to reach a particular state may be determined in advance without iterative techniques or Numerical Methods (e.g., the Euler numerical method). Given a prior voltage state v 0 , the time delay until voltage state v f is reached is given by: [0000] Δ   t = τ ρ  log  v f + q ρ v 0 + q ρ ( 13 ) [0065] If a spike is defined as occurring at the time the voltage state v reaches v S , then the closed-form solution for the amount of time, or relative delay, until a spike occurs as measured from the time that the voltage is at a given state v is: [0000] Δ   t S = { τ +  log  v S + q + v + q + if   v > v ^ + ∞ otherwise ( 14 ) [0000] where {circumflex over (v)} + is typically set to parameter v + , although other variations may be possible. [0066] The above definitions of the model dynamics depend on whether the model is in the positive or negative regime. As mentioned, the coupling and the regime p may be computed upon events. For purposes of state propagation, the regime and coupling (transformation) variables may be defined based on the state at the time of the last (prior) event. For purposes of subsequently anticipating spike output time, the regime and coupling variable may be defined based on the state at the time of the next (current) event. [0067] There are several possible implementations of the Cold model, and executing the simulation, emulation or model in time. This includes, for example, event-update, step-event update, and step-update modes. An event update is an update where states are updated based on events or “event update” (at particular moments). A step update is an update when the model is updated at intervals (e.g., 1 ms). This does not necessarily require iterative methods or Numerical methods. An event-based implementation is also possible at a limited time resolution in a step-based simulator by only updating the model if an event occurs at or between steps or by “step-event” update. Parameter Evaluation [0068] Development of neural network models may include specifications for evaluating various combinations of parameters. For example, during a design process, neural networks may be tested with various parameters to improve the design of the neural network. The parameters may include the weight of synapses, delays, neuron model parameters, parameters describing elements, such as a spike threshold and/or maximum conductance, parameters describing inputs, such as a number of images and/or audio files used, learning parameters, and/or other variables of the neural network. By setting the parameters to different values and performing testing with those different values, the neural networks may be evaluated based on various criteria, such as the efficiency of the neural network. [0069] As an example, the parameters for an exemplary neural network may be set as parameters A-D. The neural network may use a function (F) to evaluate the parameters. For example, a function F(A,B,C,D), evaluates the parameters A-D where the parameters may be given any value, such as A=1, 2, . . . 10, etc. To evaluate the neural network, it is desirable to try as many combinations of the parameters as possible. [0070] In a typical neural network, the parameter search suffers from combinatorial explosion. That is, the amount of time to evaluate a function, such as function (F), is equal to the product of the amount of time to evaluate all of the variables for a first parameter, such as parameter (A), the amount of time to evaluate all of the variables for a second parameter, such as parameter (B), and so forth. In this example, the equation for the evaluation time of the entire function (F) may be expressed as T(F)=T(A)*T(B)*T(C)*T(D), where T represents time. The parameter search may be referred to as a parameter sweep. Depending on the number of potential values for each parameter, T(F) may become very large. [0071] FIG. 5A illustrates an example of a neural network 500 . As shown in FIG. 5A , the neural network may include multiple neurons 502 that are interconnected via synapses 504 . In some cases, a first group of interconnected neurons may have a feed forward connection with a second group of interconnected neurons. That is, the second group of interconnected neurons has a one-way dependency with the first group of interconnected neurons. More specifically, the output of the first group of interconnected neurons affects the second group of interconnected neurons. Additionally, the output of the second group of interconnected neurons does not affect the first group of interconnected neurons. Thus, the first group of interconnected neurons may be designated as a first hidden sub-system of neurons and the second group of interconnected neurons may be a second hidden sub-system of neurons. [0072] FIG. 5B illustrates an example neural network that has been compartmentalized into hidden sub-systems. Specifically, the neural network of FIG. 5A has been compartmentalized into hidden sub-systems to create the exemplary neural network 530 of FIG. 5B . Based on an aspect of the present disclosure, the neural network 530 may be analyzed to determine the hidden sub-systems 506 , 508 , 510 of the neural network 530 . [0073] The three hidden sub-systems 506 , 508 , 510 may be referred to as hidden sub-systems because the sub-systems are not readily determined to the developer of the neural network and/or the untrained eye. The sub-systems may be referred to as components or restricted sets. Although FIGS. 5A and 5B only illustrate one feed forward connection 520 , 522 from each hidden sub-system 506 , 510 to the second hidden sub-system 508 , the present application is not limited to sub-systems only having one feed forward connection and is also contemplated for multiple feed forward connections. [0074] Specifically, as shown in FIG. 5B , the interconnected neurons of the first hidden sub-system 506 have a feed forward connection 520 to the interconnected neurons of the second hidden sub-system 508 . That is, the output of first hidden sub-system 506 affects the second hidden sub-system 508 . Additionally, the interconnected neurons of the third hidden sub-system 510 have a feed forward connection 522 to the interconnected neurons of the second hidden sub-system 508 . That is, the output of third hidden sub-system 510 affects the second hidden sub-system 508 . Thus, the neural network 530 may be compartmentalized to the three hidden sub-systems 506 , 508 , 510 of interconnected neurons. [0075] In the present example, each sub-system has two parameters of interest, A and B, with ten possible values (1 through 10). A single testing run to evaluate one sub-system with one of the possible values for one of the parameters may take one minute. Thus, it would take two years (i.e., 10 6 runs) to evaluate all possible combinations for all of the parameters in the neural network 530 . This is obviously impractical. Therefore, there is a need to decrease the time for evaluating all possible parameters of a neural network. [0076] According to an aspect of the present disclosure, to improve the evaluation time for possible parameters of a neural network, as shown in FIG. 5B , the network may be segregated into functionally isolated sub-networks. In one configuration, the parameters of the sub-systems may be analyzed to determine the parameters that may be pruned (e.g., removed). [0077] FIGS. 6A-6D illustrate block diagrams for segregating a neural network to functionally isolated sub-networks according to an aspect of the present disclosure. [0078] As shown in FIG. 6A , a neural network (F) 602 receives inputs for parameters A, B, C, and D. Furthermore, the neural network 602 may have an output 604 . As shown in FIG. 6B , after identifying the neural network 602 and the parameter inputs A-D, the neural network 602 may be serialized to determine the layers of the neural network 602 . The layers of the neural network are similar to the hidden sub-systems of FIG. 5B . The layers may include a first layer (e.g., input layer) and lower layers, such as a second layer, third layer, etc. [0079] The first layer refers to a group of interconnected neurons that have no internal inputs and project to another group of interconnected neurons. In FIG. 6B , the first serialized layer (S1) 606 is a first group of interconnected neurons that have no internal inputs. The first serialized layer 606 may project to a second serialized layer (S2) 608 (e.g., lower layer). The lower layer refers to a second group of interconnected neurons that receives an input from a higher layer. That is, the second serialized layer 608 is a group of interconnected neurons that receives an input from the first serialized layer 606 . Furthermore, the second serialized layer 608 does not input into the first serialized layer 606 . A neural network, such as the neural network 602 of FIG. 6 may have many serialized layers. In this example, only two serialized layers 606 and 608 are illustrated. [0080] As shown in FIG. 6C , after serializing the layers, each serialized layer 606 , 608 may be parallelized. That is, using reciprocal connections, groups or islands of neurons may be determined in each serialized layer. Specifically, for parallelization the interdependencies of a layer, or serialized layer, are determined. Based upon these interdependencies, neurons may be grouped into parallelized clusters where the output of one parallelized cluster is independent from the output of another parallelized cluster. In this manner, neurons may be grouped into parallel clusters within a serialized layer. Neurons may also be grouped into parallelized clusters without serialization of the layers. [0081] For example, as shown in FIG. 6C , the first serialized layer 606 may be parallelized into a first parallelized cluster (S11) 610 and a second parallelized cluster (S12) 612 . In this configuration, the first parallelized cluster 610 and a second parallelized cluster 612 are part of the first serialized layer 606 . Still, because the output of the neurons of the first parallelized cluster 610 is independent from the output of the neurons in the second parallelized cluster 612 , the neurons of the first serialized layer 606 may be parallelized to the aforementioned parallelized clusters. The first parallelized cluster 610 and the second parallelized cluster 612 may be referred to as a parallelized block P1) 614 . [0082] The serialized layers and parallelized clusters may be referred to as sub-systems. In one configuration, once the sub-systems of the neural network have been determined, the parameters of the neural network are pruned. That is, the parameters of each sub-system are analyzed to determine if the parameters affect the output of the sub-system. If a parameter does not affect a sub-system, the parameter may be pruned for purposes of evaluating the respective sub-system as that parameter's value does not alter the eventual output of the respective sub-system. [0083] As shown in FIG. 6D , the inputs to the first parallelized cluster 610 , second parallelized cluster 612 , and second serialized layer 608 (e.g., the sub-systems) are pruned. In this example, the system may determine that parameter A only affects the first parallelized cluster 610 , parameters B and D only affect the second parallelized cluster 612 , and parameters C and D only affect the second serialized layer 608 . Thus, the parameters not affecting a specific sub-system are pruned and removed as inputs to the respective sub-system. [0084] In one configuration, the identification of layers may be recursively performed. That is, each serialized layer may be treated as a serialized input layer to determine other serialized layers. Specifically, as previously discussed, the first layer refers to a group of interconnected neurons that have no internal inputs. The first layer also outputs to another group of interconnected neurons (e.g., second layer of interconnected neurons). That is, the first layer of neurons is not affected by the outputs of any of the other neurons in the network. Moreover, the first layer of neurons may receive inputs from outside of the neural network. [0085] For example, a user may input image files that are transduced to a neural spiking representation to the first layer. After identifying the first layer of neurons, the second layer of neurons that receive inputs from neurons in the first layer are identified. The second layer of neurons only receives inputs from the first layer of neurons. That is, the inputs of the first layer of neurons are treated as external inputs. Thus, the search for the different layer of neurons is recursively performed based on the layers that only receive inputs from previous layers and have no other internal inputs. [0086] The present disclosure is not limited to both serializing layers and parallelizing layers. In one configuration, the neural network is only serialized and pruned. Alternatively, in another configuration, the neural network is only parallelized and pruned. In another configuration, the user may identify the sub-systems. In still yet another configuration, the identification of the sub-systems may be identified by the compiler when the neural network is compiled. [0087] As described above, the sub-system may be defined based on network connectivity. Additionally or alternatively, in one configuration, the sub-system may be defined by monitoring activity levels of a sweep. That is, when a first layer sub-system receives an input from a second layer sub-system, the first layer sub-system computes its outputs over all the stored outputs of the lower layer. If some of the outputs of the first layer (e.g., inputs to the second layer) are similar to each other (e.g. no spiking output for specific parameter ranges) then the system may run the computation once for all degenerate inputs. Moreover, the output of the second layer for one run may be stored as the output of the combined run for the entire range that produced the degenerate output from the first layer. [0088] Furthermore, in one configuration, after each step in a sweep, an output of each sub-system is stored. That is, the parameter sweep may to produce a large collection of outputs of the combined system. Accordingly, the user may desire to analyze the outputs according to a function that is defined on the outputs of the combined network. For example, the user may desire to analyze whether neurons respond selectively to a pattern of spiking outputs of individual neurons in the final serialized layer. [0089] FIG. 7 illustrates a method 700 for improving a parameter evaluation of a neural network according to an aspect of the present disclosure. At block 702 , the method serializes the neural network into layers. In one configuration, the serialization determines a first layer (e.g., input layer) that receives an input of the parameters but has no other internal inputs. Additionally, in this configuration, the serialization of block 702 also determines a second layer that receives an input from the first layer. After serializing a layer, the serialized layer is parallelized at block 704 . In one configuration, the parallelization uses reciprocal connections to determine clusters of neurons within the serialized layer. Each cluster of neurons within the serialized layer may be referred to as a parallelized cluster. [0090] In block 706 , the parameters of the neural network are pruned. The pruning refers to removing parameters that are not used by a specific parallelized cluster and/or serialized layer. As previously discussed the parallelized clusters and serialized layers of a neural network may be referred to as sub-systems, restricted sets, or components. In block 708 , the computation time (Te) and storage requirements (S, Ta) for each sub-system are determined. In one configuration, after determining the computation time (Te) and storage requirements (S, Ta) for each sub-system at block 708 , blocks 702 - 706 may be performed recursively to determine additional sets (parallelized and/or serialized). In another configuration, at block 710 , an objective function determines the sub-systems to cache based on access time, memory access time, and/or computational time. [0091] FIG. 8 illustrates functions for evaluating the parameters based on an aspect of the present disclosure. The flow chart of FIG. 8 is based on the block diagrams of FIGS. 6A-D . As shown in block 802 , the parameters (A, B, C, and D) of a neural network may be evaluated based on a function F(A,B,C,D), where F is the function to evaluate the parameters. The output of the first serialized layer (S1) 606 is a function of its inputs, parameters A, B, C, and D and may thus be expressed as S1(A,B,C,D). The output of the second serialized layer 608 may also be expressed as a function of its inputs. As the inputs to the second serialized layer are parameters A, B, C, and D, in addition to the output from the first serialized layer 606 , the output of the second serialized layer may be expressed as S2(S1(A,B,C,D), A,B,C,D). As shown in block 802 , because the second serialized layer 608 is the final layer of the system, the function F(A,B,C,D) is equal to the output of the second serialized layer 608 . Therefore, the output of the system function F(A,B,C,D) may be said to be equal to S2(S1(A,B,C,D), A,B,C,D). [0092] At block 804 , the function for the first serialized layer 606 is determined. That is, the first serialized layer 606 has a function of S1(A,B,C,D) which is equal to functions of a parallelized block 614 (P1). The parallelized block 614 includes the first parallelized cluster 610 (S11) and the second parallelized cluster 612 (S12). As shown in FIGS. 6A-D , the first parallelized cluster 610 and the second parallelized cluster 612 each have inputs of the parameters A, B, C, and D. Therefore, the first serialized layer 606 may be equal to a function P1(S11(A,B,C,D), S12(A,B,C,D)). [0093] At block 806 , pruning is performed to improve the evaluation of the parameters. The first parallelized cluster (S11) 610 may be pruned to only be a function of parameter A, therefore, the first parallelized cluster 610 is equal to S11(A). Additionally, the second parallelized cluster (S12) 612 may be pruned to only be a function of parameters B and D. Therefore, the second parallelized cluster 612 is equal to S12(B,D). Additionally, the second serialized layer 608 is pruned so that the inputs are only parameters C and D. [0094] After pruning, the function for the first parallelized block 620 is P1((S11(A), S12(B,D)). Likewise, the function for the second serialized layer is S2(P1(S11(A), S12(B,D)),C,D). After pruning, at block 808 the evaluation times (Te) for the layers and the neural network are determined. That is, the evaluation time for the first serialized layer 606 (S1) is the sum of the evaluation of the first parallelized cluster 610 and the second parallelized cluster 612 (S12). Furthermore, the neural network (F) is the product of the evaluation time for the first serialized layer 606 (S1) and the second serialized layer 608 (S2). [0095] In one configuration, at block 810 an optimization function is executed to determine the sub-systems to cache based on the access time, memory access time, and/or computational time. The optimization function can be based on the following equation: [0000] O ( N )=alpha*( Te−Ta )− S   (15) [0096] In equation (15), 0 is the optimization function and N is the identified sub-systems, such as a first serialized layer 606 . Te is an evaluation time for input parameters. Ta is an evaluation time for parameters access from memory. [0097] In equation (15) S is a measure of an amount of memory storage specified for each subsystem. For each sub-system, Te is the execution time of that subsystem. The execution time includes accessing the inputs of the subsystem in addition to generating and storing the outputs. Ta is the time to access the outputs of the system if they were stored in memory rather than recomputed from their inputs. If memory space is not a concern, the system may still compute memory access times to determine whether the system should cache the outputs from a particular subsystem. [0098] In one configuration, the pruning may be performed to improve the computation along different dimensions because there may be different costs associated with computation time or memory storage. For example, a user may select to re-compute the outputs of a sub-system rather than store the outputs if the user is allocated more computing power and less memory storage. In this example, the trade-off between computation power and memory storage may be set as “alpha” in equation (15). In another configuration, the system determines to cache the outputs of subsystems for which the time to access the outputs Ta is less than the time to re-compute the outputs Te. In this configuration, the default settings in equation (15) are alpha is equal to 1 and S is equal to 0. [0099] FIG. 9 illustrates an example implementation 900 of the aforementioned modifying neural dynamics using a general-purpose processor 902 in accordance with certain aspects of the present disclosure. Variables (neural signals), synaptic weights, system parameters associated with a computational network (neural network), delays, and/or frequency bin information may be stored in a memory block 904 , while instructions executed at the general-purpose processor 902 may be loaded from a program memory 909 . In an aspect of the present disclosure, the instructions loaded into the general-purpose processor 902 may comprise code for obtaining prototypical neuron dynamics and/or modifying parameters of a neuron model so that the neuron model matches the prototypical neuron dynamics. [0100] FIG. 10 illustrates an example implementation 1000 of the aforementioned modifying neural dynamics where a memory 1002 can be interfaced via an interconnection network 1004 with individual (distributed) processing units (neural processors) 1006 of a computational network (neural network) in accordance with certain aspects of the present disclosure. Variables (neural signals), synaptic weights, system parameters associated with the computational network (neural network) delays, and/or frequency bin information may be stored in the memory 1002 , and may be loaded from the memory 1002 via connection(s) of the interconnection network 1004 into each processing unit (neural processor) 1006 . In an aspect of the present disclosure, the processing unit 1006 may be configured to obtain prototypical neuron dynamics and/or modify parameters of a neuron model. [0101] FIG. 11 illustrates an example implementation 1100 of the aforementioned modifying neural dynamics. As illustrated in FIG. 11 , one memory bank 1102 may be directly interfaced with one processing unit 1104 of a computational network (neural network). Each memory bank 1102 may store variables (neural signals), synaptic weights, and/or system parameters associated with a corresponding processing unit (neural processor) 1104 delays, and frequency bin information. In an aspect of the present disclosure, the processing unit 1104 may be configured to obtain prototypical neuron dynamics and/or modify parameters of a neuron model. [0102] FIG. 12 illustrates an example implementation of a neural network 1200 in accordance with certain aspects of the present disclosure. As illustrated in FIG. 12 , the neural network 1200 may have multiple local processing units 1202 that may perform various operations of the methods described above. Each processing unit 1202 may be a local state memory 1204 and a local parameter memory 1206 that store states and parameters of the neural network. In addition, the processing unit 1202 may have a memory 1208 that stores a local (neuron) model program, a memory 1210 that stores a local learning program, and a local connection memory 1212 . Furthermore, as illustrated in FIG. 12 , each local processing unit 1202 may be interfaced with a unit 1214 for configuration processing that may provide configuration for local memories of the local processing unit, and with routing connection processing elements 1216 that provide routing between the local processing units 1202 . [0103] According to certain aspects of the present disclosure, each local processing unit 1202 may be configured to determine parameters of the neural network based upon desired one or more functional features of the neural network. Each local processing unit may develop the one or more functional features towards the desired functional features as the determined parameters are further adapted, tuned, and updated. [0104] FIG. 13 illustrates a method 1300 for performing a parameter sweep over a system having sub-systems with one-way dependencies. In one configuration, at block 1302 the neuron model serializes sub-systems of the system by determining the one way dependencies between the sub-systems. Furthermore, in addition to or alternate from the serializing of block 1302 , at block 1304 , the neuron model parallelizes the sub-systems by determining independencies within each sub-system. That is, in one configuration, the neuron model may either perform the serializing of block 1302 or the parallelizing of block 1304 . In another configuration, the neuron model may perform both the serializing of block 1302 and the parallelizing of block 1304 . Finally, at block 1306 , the neuron model prunes input parameters of each sub-system based on whether each input parameter affects each sub-system. [0105] In one configuration, a neural network, such as the neural network of the aspects of the present disclosure, is configured to modify neural dynamics. The neural network may include means for parallelizing, means for serializing, and/or means for pruning. In one aspect, the parallelizing means, serializing means, and/or pruning means may be the program memory 906 , memory block 1204 , memory 1002 , interconnection network 1004 , processing units 1006 , processing unit 1104 , local processing units 1202 , and or the routing connection processing elements 916 configured to perform the functions recited by the parallelizing means, serializing means, and/or pruning means [0106] The various operations of methods described above may be performed by any suitable means capable of performing the corresponding functions. The means may include various hardware and/or software component(s) and/or module(s), including, but not limited to, a circuit, an application specific integrated circuit (ASIC), or processor. Generally, where there are operations illustrated in Figures, those operations may have corresponding counterpart means-plus-function components with similar numbering. [0107] As used herein, the term “determining” encompasses a wide variety of actions. For example, “determining” may include calculating, computing, processing, deriving, investigating, looking up (e.g., looking up in a table, a database or another data structure), ascertaining and the like. Additionally, “determining” may include receiving (e.g., receiving information), accessing (e.g., accessing data in a memory) and the like. Furthermore, “determining” may include resolving, selecting, choosing, establishing and the like. [0108] As used herein, a phrase referring to “at least one of” a list of items refers to any combination of those items, including single members. As an example, “at least one of: a, b, or c” is intended to cover: a, b, c, a-b, a-c, b-c, and a-b-c. [0109] The various illustrative logical blocks, modules and circuits described in connection with the present disclosure may be implemented or performed with a general purpose processor, a digital signal processor (DSP), an application specific integrated circuit (ASIC), a field programmable gate array signal (FPGA) or other programmable logic device (PLD), discrete gate or transistor logic, discrete hardware components or any combination thereof designed to perform the functions described herein. A general-purpose processor may be a microprocessor, but in the alternative, the processor may be any commercially available processor, controller, microcontroller or state machine. A processor may also be implemented as a combination of computing devices, e.g., a combination of a DSP and a microprocessor, a plurality of microprocessors, one or more microprocessors in conjunction with a DSP core, or any other such configuration. [0110] The steps of a method or algorithm described in connection with the present disclosure may be embodied directly in hardware, in a software module executed by a processor, or in a combination of the two. A software module may reside in any form of storage medium that is known in the art. Some examples of storage media that may be used include random access memory (RAM), read only memory (ROM), flash memory, EPROM memory, EEPROM memory, registers, a hard disk, a removable disk, a CD-ROM and so forth. A software module may comprise a single instruction, or many instructions, and may be distributed over several different code segments, among different programs, and across multiple storage media. A storage medium may be coupled to a processor such that the processor can read information from, and write information to, the storage medium. In the alternative, the storage medium may be integral to the processor. [0111] The methods disclosed herein comprise one or more steps or actions for achieving the described method. The method steps and/or actions may be interchanged with one another without departing from the scope of the claims. In other words, unless a specific order of steps or actions is specified, the order and/or use of specific steps and/or actions may be modified without departing from the scope of the claims. [0112] The functions described may be implemented in hardware, software, firmware, or any combination thereof. If implemented in hardware, an example hardware configuration may comprise a processing system in a device. The processing system may be implemented with a bus architecture. The bus may include any number of interconnecting buses and bridges depending on the specific application of the processing system and the overall design constraints. The bus may link together various circuits including a processor, machine-readable media, and a bus interface. The bus interface may be used to connect a network adapter, among other things, to the processing system via the bus. The network adapter may be used to implement signal processing functions. For certain aspects, a user interface (e.g., keypad, display, mouse, joystick, etc.) may also be connected to the bus. The bus may also link various other circuits such as timing sources, peripherals, voltage regulators, power management circuits, and the like, which are well known in the art, and therefore, will not be described any further. [0113] The processor may be responsible for managing the bus and general processing, including the execution of software stored on the machine-readable media. The processor may be implemented with one or more general-purpose and/or special-purpose processors. Examples include microprocessors, microcontrollers, DSP processors, and other circuitry that can execute software. Software shall be construed broadly to mean instructions, data, or any combination thereof, whether referred to as software, firmware, middleware, microcode, hardware description language, or otherwise. Machine-readable media may include, by way of example, RAM (Random Access Memory), flash memory, ROM (Read Only Memory), PROM (Programmable Read-Only Memory), EPROM (Erasable Programmable Read-Only Memory), EEPROM (Electrically Erasable Programmable Read-Only Memory), registers, magnetic disks, optical disks, hard drives, or any other suitable storage medium, or any combination thereof. The machine-readable media may be embodied in a computer-program product. The computer-program product may comprise packaging materials. [0114] In a hardware implementation, the machine-readable media may be part of the processing system separate from the processor. However, as those skilled in the art will readily appreciate, the machine-readable media, or any portion thereof, may be external to the processing system. By way of example, the machine-readable media may include a transmission line, a carrier wave modulated by data, and/or a computer product separate from the device, all which may be accessed by the processor through the bus interface. Alternatively, or in addition, the machine-readable media, or any portion thereof, may be integrated into the processor, such as the case may be with cache and/or general register files. Although the various components discussed may be described as having a specific location, such as a local component, they may also be configured in various ways, such as certain components being configured as part of a distributed computing system. [0115] The processing system may be configured as a general-purpose processing system with one or more microprocessors providing the processor functionality and external memory providing at least a portion of the machine-readable media, all linked together with other supporting circuitry through an external bus architecture. Alternatively, the processing system may comprise one or more neuromorphic processors for implementing the neuron models and models of neural systems described herein. As another alternative, the processing system may be implemented with an ASIC (Application Specific Integrated Circuit) with the processor, the bus interface, the user interface, supporting circuitry, and at least a portion of the machine-readable media integrated into a single chip, or with one or more FPGAs (Field Programmable Gate Arrays), PLDs (Programmable Logic Devices), controllers, state machines, gated logic, discrete hardware components, or any other suitable circuitry, or any combination of circuits that can perform the various functionality described throughout this disclosure. Those skilled in the art will recognize how best to implement the described functionality for the processing system depending on the particular application and the overall design constraints imposed on the overall system. [0116] The machine-readable media may comprise a number of software modules. The software modules include instructions that, when executed by the processor, cause the processing system to perform various functions. The software modules may include a transmission module and a receiving module. Each software module may reside in a single storage device or be distributed across multiple storage devices. By way of example, a software module may be loaded into RAM from a hard drive when a triggering event occurs. During execution of the software module, the processor may load some of the instructions into cache to increase access speed. One or more cache lines may then be loaded into a general register file for execution by the processor. When referring to the functionality of a software module below, it will be understood that such functionality is implemented by the processor when executing instructions from that software module. [0117] If implemented in software, the functions may be stored or transmitted over as one or more instructions or code on a computer-readable medium. Computer-readable media include both computer storage media and communication media including any medium that facilitates transfer of a computer program from one place to another. A storage medium may be any available medium that can be accessed by a computer. By way of example, and not limitation, such computer-readable media can comprise RAM, ROM, EEPROM, CD-ROM or other optical disk storage, magnetic disk storage or other magnetic storage devices, or any other medium that can be used to carry or store desired program code in the form of instructions or data structures and that can be accessed by a computer. In addition, any connection is properly termed a computer-readable medium. For example, if the software is transmitted from a website, server, or other remote source using a coaxial cable, fiber optic cable, twisted pair, digital subscriber line (DSL), or wireless technologies such as infrared (IR), radio, and microwave, then the coaxial cable, fiber optic cable, twisted pair, DSL, or wireless technologies such as infrared, radio, and microwave are included in the definition of medium. Disk and disc, as used herein, include compact disc (CD), laser disc, optical disc, digital versatile disc (DVD), floppy disk, and Blu-ray® disc where disks usually reproduce data magnetically, while discs reproduce data optically with lasers. Thus, in some aspects computer-readable media may comprise non-transitory computer-readable media (e.g., tangible media). In addition, for other aspects computer-readable media may comprise transitory computer-readable media (e.g., a signal). Combinations of the above should also be included within the scope of computer-readable media. [0118] Thus, certain aspects may comprise a computer program product for performing the operations presented herein. For example, such a computer program product may comprise a computer-readable medium having instructions stored (and/or encoded) thereon, the instructions being executable by one or more processors to perform the operations described herein. For certain aspects, the computer program product may include packaging material. [0119] Further, it should be appreciated that modules and/or other appropriate means for performing the methods and techniques described herein can be downloaded and/or otherwise obtained by a user terminal and/or base station as applicable. For example, such a device can be coupled to a server to facilitate the transfer of means for performing the methods described herein. Alternatively, various methods described herein can be provided via storage means (e.g., RAM, ROM, a physical storage medium such as a compact disc (CD) or floppy disk, etc.), such that a user terminal and/or base station can obtain the various methods upon coupling or providing the storage means to the device. Moreover, any other suitable technique for providing the methods and techniques described herein to a device can be utilized. [0120] It is to be understood that the claims are not limited to the precise configuration and components illustrated above. Various modifications, changes and variations may be made in the arrangement, operation and details of the methods and apparatus described above without departing from the scope of the claims.

An artificial neural network may be configured to test the impact of certain input parameters. To improve testing efficiency and to avoid test runs that may not alter system performance, the effect of input parameters on neurons or groups of neurons may be determined to classify the neurons into groups based on the impact of certain parameters on those groups. Groups may be ordered serially and/or in parallel based on the interconnected nature of the groups and whether the output of neurons in one group may affect the operation of another. Parameters not affecting group performance may be pruned as inputs to that particular group prior to running system tests, thereby conserving processing resources during testing.

CROSS REFERENCE TO RELATED APPLICATION [0001] This application claims the benefit of U.S. Provisional Patent Application No. 60/743,456, filed Mar. 10, 2006, the disclosure of which is incorporated by reference herein in its entirety. TECHNICAL FIELD [0002] This invention relates to a method and a device for cleaning griddles using disposable scrubbing pads. BACKGROUND OF THE INVENTION [0003] Restaurants commonly have one or more griddle surfaces that provide a flat, hot cooking surface for cooking food items. Often restaurants include both a flat griddle to cook foods such as eggs and pancakes and a grooved griddle to cook meats where a charbroiled appearance is desirable. In addition to the aesthetic appeal associated with food cooked over a grooved griddle, the grooved griddle is preferable over a flat griddle when draining fat out of meat products while cooking the meat is desirable. When cooking meat products on a grooved griddle the meat product rests over raised ridges on the griddle. As the meat product cooks, the fat drains from the meat and collects on the lower surfaces of the griddle that are positioned between the raised ridges on the griddle. Though traditional open flame grills also enable fat to drain from meat products while the meat is cooking, grooved griddles are sometime preferred over traditional open flame grills because they are typically more energy efficient and the temperature of the cooking surface can be more easily controlled. [0004] Cleaning tools have been developed to remove the buildup of grease and food particles on griddles and open flame grills. Exemplary tools are disclosed in U.S. Pat. No. 6,966,094 to Rigakos; U.S. Pat. No. 6,871,377 to Veltrop et al; U.S. Pat. No. 6,443,646 to MacDonald; U.S. Pat. No. 6,351,887 to Hurst; U.S. Pat. No. 6,263,578 to Frantz et al.; U.S. Pat. No. 6,216,306 to Esterson et al.; U.S. Pat. No. 6,039,372 to Noe et al.; U.S. Pat. No. 5,373,600 to Stojanovski et al.; U.S. Pat. No. 5,255,406 to Rood; U.S. Pat. No. 4,668,302 to Kolodziej et al.; U.S. Pat. No. 4,516,870 to Nakozato; U.S. Pat. No. 4,146,943 to Werthermer et al.; U.S. Pat. No. 4,071,983 to Thielen; U.S. Pat. No. 4,056,863 to Gunjian; and U.S. Pat. No. D470,985 to Zemel. Known tools are not particularly well suited for cleaning grooved griddles of various geometric configurations. [0005] Grooved griddles are difficult to clean with tools designed to clean flat griddles or grills. Typically, such tools have problems cleaning the area between the raised portions of the griddle. Known tools for cleaning grooved griddles are less than effective because griddles are not uniform in size or geometric configuration. In addition, known tools often require the user to be positioned too close to the hot griddle surface. Moreover, the useful life and versatility of the entire tool is typically limited by the cleaning element of the tool. Accordingly, there is a need for improved cleaning devices that enable a user to clean a grooved griddle more efficiently and effectively. SUMMARY OF THE INVENTION [0006] The invention provides a cleaning element configured to attach to an end of a griddle cleaning tool. The cleaning elements according to the invention are configured to efficiently and effectively clean an uneven grooved griddle surface. BRIEF DESCRIPTION OF THE DRAWINGS [0007] FIG. 1 is a perspective assembly view of a griddle cleaning tool including a pad according to an embodiment of the invention positioned over a grooved griddle; [0008] FIG. 2 is an end view of the pad shown in FIG. 1 ; [0009] FIG. 3 a is a perspective view of an alternative embodiment of the pad shown in FIG. 1 ; [0010] FIG. 3 b is a side elevation view of a portion of a grooved griddle surface; [0011] FIG. 3 c is side elevation view of a portion of a grooved griddle surface; [0012] FIG. 4 a is an end view of the pad shown in FIG. 3 a; [0013] FIG. 4 b is side elevation view of the pad shown in FIG. 3 a on the grooved griddle surface shown in FIG. 3 b; [0014] FIG. 4 c is side elevation view of the pad shown in FIG. 3 a on the grooved griddle surface shown in FIG. 3 b; [0015] FIG. 5 is a perspective view of another alternative embodiment of the pad shown in FIG. 1 ; [0016] FIG. 6 is an end view of the pad shown in FIG. 5 ; [0017] FIG. 7 is a perspective view of an alternative embodiment of a griddle cleaning tool shown in FIG. 1 ; [0018] FIG. 8 is a perspective view of an alternative embodiment of a griddle cleaning tool shown in FIG. 1 ; [0019] FIG. 9 is a perspective assembly view of the pad shown in FIG. 8 ; [0020] FIG. 10 is a top perspective view of a portion of the griddle cleaning tool in FIG. 8 ; [0021] FIG. 11 is a bottom perspective view of a portion of the griddle cleaning tool in FIG. 8 ; and [0022] FIG. 12 is a top perspective view of an alternative embodiment of the portion of the griddle cleaning tool in FIG. 10 . DETAILED DESCRIPTION [0023] Referring to FIG. 1 , a griddle cleaning tool 10 is shown. The tool includes a handle 12 and a foot 14 . The bottom surface 16 of the foot 14 includes a plurality of hooks 18 , which are configured to engage and secure the pad 20 on the bottom surface 16 of the foot 14 . In the embodiment shown the foot 14 and the handle 12 is one piece. In some embodiments the handle and the foot are separate pieces. See Application Ser. No. 60/743,455 docket number 61852US002 having the same filing date as this application in the name of 3M Innovative Properties Company, the subject matter of which is incorporated herein by reference. [0024] Referring to FIGS. 2 and 3 a - c , the pad 20 shown in FIG. 1 is generally rectangular in shape and includes a stepped cross sectional profile. The pad 20 includes peaks 22 separated by valleys 24 . The peaks 22 and valleys 24 of pad 20 include flat top surfaces 26 and 28 and are spaced apart by a distance D 1 . Preferred distance D 1 is constant across the pad 20 and matches the griddle groove spacing GGS of whatever griddle model that the pad 20 is designed to clean. Preferably, the spacing D 1 is within +/−20% of the groove spacing GGS. Since not all griddles have the same griddle spacing GGS, the pad 20 can be manufactured in several sizes with various peak and valley spacing to accommodate particular differences in griddle spacing. When the pad 20 is in use, the peaks 22 of the pad 20 contacts the low portions 30 (shown in FIG. 1 ) of the griddle 34 and the valleys 24 of the pad 20 engage the high portions 32 (shown in FIG. 1 ) of the griddle 34 . [0025] Still referring to FIGS. 2 and 3 a - c , the top surface (commonly referred to as the back surface) 36 ( 58 and 68 in FIGS. 3 b and 3 c respectively) of the pad 20 is configured to be secured to the bottom surface 16 of the foot 14 via the plurality of hooks 18 . Once the pad 20 is secured to the foot 14 , the griddle can be cleaned by moving the handle 12 back and forth across the griddle 34 until the pad 20 breaks loose the food, grease, and carbonized material from the griddle 34 surface. [0026] It should be understood that the hooks 18 of the foot 14 need not be in the shape shown in the figures, but that the hooks 18 can be in any geometric configuration capable of engaging and securing the pad 20 to the foot 14 . In addition, in alternative embodiments the foot 14 may have no hooks 18 . Instead, the pad may include an adhesive strip or other engagement mechanisms that secure the pad 20 to the foot or it may include clamps for securing the edge of the pad 20 to the foot 14 . [0027] In some embodiments the pad 20 comprises a non-woven substrate suited for scouring heated surfaces. In some embodiments the non-woven substrate also includes solid cleaners disposed therein or thereon that at least partially remove or soften the food soils. In many embodiments, non-woven substrates include non-woven webs of fibers. [0028] In some embodiments the pad 20 can be used in conjunction with a liquid or a solid chemical cleaner. For example, the pad 20 can be used with 3M's commercially available Scotch-Brite Quick Clean Griddle Liquid, which is griddle cleaning liquid intended for use on food contact surfaces and is useful in loosening and lifting carbonized grease and food soil from hot griddle surfaces. In other embodiments, the pad 20 can be impregnated or otherwise attached to a chemical cleaner. [0029] In one embodiment the pad 20 includes features disclosed in PCT Publication Number WO 2007/101866 (3M Innovative Properties Company). The entire PCT filing is incorporated by reference herein and portions of the application are included below. [0030] The following disclosure is believed to be applicable generally to solid cleaners and the use of such solid cleaners on heated surfaces. Specifically, the disclosure is based around a solid cleaner that melts on a heated food preparation surface such as, for example, a grill surface, a griddle surface, or an oven surface. The heated surface can be formed of any material including, for example, metal, ceramic, glass, and/or plastic. These examples, and the examples discussed below, provide an appreciation of the applicability of the disclosed cleaning systems, but should not be interpreted in a limiting sense. [0031] A solid cleaner for heated surfaces is disclosed that includes one or more D solidifying agents and one or more cleaning agents. The solid cleaner is solid at room temperature (e.g., 24 degrees Celsius) and a liquid at an elevated temperature. The elevated temperature can be any useful temperature at which the solid cleaner begins to melt (e.g., melting point.) The solid cleaner can have any useful melting point. In some embodiments, the solid cleaner has a melting point in a range from 35 to 150 degrees Celsius or from 35 to 100 degrees Celsius, or from 45 to 90 degrees Celsius, as desired. Solid cleaners that melt on heated surfaces provide one or more of the following advantages over liquid cleaners: increased dwell time; decreased cleaner evaporation; and/or the ability to be used on vertical heated surfaces. In many embodiments, the solid cleaners have an accelerated cleaning action at elevated temperatures (e.g., above 100 degrees Celsius). In many embodiments, the solid cleaner is generally recognized as safe (GRAS) for food contact. [0032] The solid cleaner can be any defined size or shape. In some embodiments, the solid cleaner has a cube shape, a cuboid shape, a pyramid shape, a cylinder shape, a cone shape, a sphere shape, or portions thereof. In some embodiments, the solid cleaner has a weight from 1 gram to 10 kilograms, or from 1 to 1000 grams, or from 5 to 500 grams, or from 10 to 200 grams. In other embodiments, the solid cleaner is a powder, pellet, flake, tablet, bar, and the like. The solid cleaner can be combined, or used in conjunction with other cleaning articles such as, for example a non-woven scouring pad, as described below, an abrasive coated woven web substrate griddle screen such as, for example SCOTCH-BRITE™ griddle screen number 200 , or a pumice block, as desired. [0033] The solid cleaner includes one or more solidifying agents that can assist in forming the solid cleaner. The term “solid” can be defined as a material having a definite volume and configuration independent of its container. Any useful solidifying agent can be used to form the solid cleaner. Any useful amount of solidifying agent can be used to assist in solidifying the solid cleaner. In many embodiments, the solidifying agent is inert or does not assist in the cleaning action of the solid cleaner. In many embodiments, the solidifying agent is generally recognized as safe (GRAS) for food contact. In certain embodiments, the solid cleaner does not need to be rinsed off of the cleaned surface, implying that it is a “no-rinse” cleaner and GRAS for food contact. [0034] In many embodiments, the solidifying agent includes one or more waxes. The wax can be a natural wax or synthetic wax. In some embodiments where the solid cleaner includes wax, the solid cleaner is substantially insoluble in water up to at least 35 degrees Celsius. In some embodiments, the solidifying agent includes a natural wax such as, for example, a beeswax, a candelilla wax, a carnauba wax, a rice bran wax, a lemon peel wax, a soy wax, an orange peel wax, or mixtures thereof. In other embodiments, the solidifying agent includes a synthetic wax such as, for example, Baker-Hugnes (Petrolite) makes Bareco High Melt Microcrystalline waxes (melting point 82 to 93 degrees Celsius), Bareco Flexible Microcrystalline waxes (melting point 65 to 82 degrees Celsius), Starwax™, Victory™, Ultraflex™ and Be Square™ waxes, among others. EMS-Griltech (Switzerland) also makes synthetic low melting polymers such as copolyamide, and copolyesters. Synthetic waxes can also include PEG waxes that are solids such as PEG 1000 NF/FCC, fatty alcohols such as cetyl alcohol, and fatty esters such as propylene glycol monostearate, glycerol monolaurate, and sorbitan esters. [0035] In some embodiments, the solidifying agent includes an emulsifying wax. The emulsifying wax can replace a portion of the one or more waxes, as desired. Emulsifying wax can include, for example, a blend of fatty acids (stearic, palmitic, oleic, capric, caprylic, myristic, and lauric), fatty alcohols (stearyl, cetyl) and/or fatty esters (polysorbates or TWEEN), and the like. In some embodiments, the emulsifying wax is a fatty alcohol such as, for example, stearic alcohol, cetyl alcohol, or mixtures thereof. One example of an emulsifying wax is Emulsifying Wax NF (cas# 67762-27-0; 9005-67-8) and is a blend of cetearyl alcohol, polysorbate 60 , PEG-150 stearate & steareth-20. If present, the emulsifying wax to other wax weight ratio can be from 1:1 to 1:5, or from 3:1 to 1:3, or from 2:1 to 1:2 as desired. [0036] Wax can be included in the solid cleaner in any useful amount. In many embodiments, a solidifying amount of wax is included in the solid cleaner. In some embodiments, wax is present in the solid cleaner in a range from 10 to 80 wt %, or from 25 to 75 wt %, or from 30 to 50 wt %. [0037] In some embodiments, the solidifying agent includes a one or more solid polyols. The term “polyol” refers to any organic molecule comprising at least two free hydroxyl groups. Polyols include polyoxyethylene derivatives such as, for example, glycol (diols), triols and monoalcohols, ester, or ethers thereof. Examples of polyols include solids glycols such as, for example, polyethylene glycols (PEG) under the tradename Carbowax series available from Dow Chemical, Midland Mich., polypropylene glycols (PPG) available from Dow Chemical, Midland, Mich., sorbitol and sugars, and solid polyesters such as, for example, poly(ε-caprolactone) under the tradename TONE series from Dow Chemical, Midland Mich., glycerol esters such as, for example, fatty acid mono ester. Fatty acid monoesters include but are not limited to propylene glycol monostearate, glycerol monolaurate, and glycerol monostearate. These esters are GRAS or approved as direct food additives. [0038] Polyol can be included in the solid cleaner in any useful amount. In many embodiments, a solidifying amount of polyol is included in the solid cleaner. In some embodiments, polyol is present in the solid cleaner in a range from 10 to 80 wt %, or from 25 to 75 wt %, or from 30 to 50 wt %. [0039] The solid cleaner includes one or more cleaning agents that can assist in the cleaning action of the solid cleaner. The cleaning agent can be any useful cleaning agent. The cleaning agent can be present in the solid cleaner in any useful amount. In many embodiments, the cleaning agents are generally recognized as safe (GRAS) for food contact. [0040] Cleaning agents include, for example, surfactants, and pH modifiers. In many embodiments, a cleaning amount of cleaning agent is included in the solid cleaner. In many embodiments, the cleaning agent is capable of removing at least a portion of the soil or residue on the heating surface without mechanical scrubbing action. In illustrative embodiments, the cleaning agent is present in the solid cleaner in range from 1 to 90 wt %, or from 1 to 50 wt %, or from 5 to 30 wt %. [0041] In some embodiments, the cleaning agent includes one or more pH modifiers. These pH modifiers include alkaline compounds such as, inorganic alkaline compounds including for example, hydroxides, silicates, phosphates, and carbonates; and organic alkaline compounds including for example, amines. In other embodiments, the pH modifier is an acidic compound such as, for example, citric acid and the like. [0042] In some embodiments, the cleaning agent is a carbonate salt such as, for example, calcium carbonate, potassium carbonate, or sodium carbonate. In some embodiments, the carbonate salt includes potassium carbonate and sodium carbonate that is dissolved in water, forming carbonate ions. In other embodiments, the carbonate salt includes a bicarbonate salt such as, for example, sodium bicarbonate. In further embodiments, the cleaning agent includes a silicate salt such as, for example, sodium metasilicate. [0043] The pH modifiers can be included in the solid cleaner in any useful amount. In many embodiments, the pH modifier is present in the solid cleaner in range from 0.1 to 80 wt %, or from 1 to 50 wt %, or from 5 to 30 wt %. In many embodiments, the solid cleaner has a pH in a range from 7 to 13. [0044] In some embodiments, the cleaning agent includes one or more surfactants. These surfactants include, for example, natural surfactants, anionic surfactants, nonionic surfactants, and amphoteric surfactants. Natural surfactants include, but are not limited to, coconut-based soap solutions. Anionic surfactants include, but are not limited to, dodecyl benzene sulfonic acid and its salts, alkyl ether sulfates and salts thereof, olefin sulfonates, phosphate esters, soaps, sulfosuccinates, and alkylaryl sulfonates. Amphoteric surfactants include, but are not limited to, imidazoline derivatives, betaines, and amine oxides. These surfactants can be included in the solid cleaner in any useful amount. In many embodiments, the surfactant is present in the solid cleaner in range from 5 to 80 wt %, or from 5 to 50 wt %, or from 5 to 30 wt %. In many embodiments, the surfactant is food grade surfactant, approved for use as a direct food additive. Often, food grade surfactants are used so that the cleaning surface does not need to be rinsed. [0045] In some embodiments, the cleaning agent includes carbonate salts such as, for example, sodium and/or potassium carbonate with an amount of surfactant less than 5 wt %, or less than 3 wt %, or less than 1 wt % based on the solid cleaner weight. In some embodiments, the cleaning agent includes carbonate salts such as, for example, sodium and/or potassium carbonate with an amount of a natural surfactant less than 5 wt %, or less than 3 wt %, or less than 1 wt % based on the solid cleaner weight. [0046] The solid cleaner may optionally include one or more carriers. The carrier can be any amount of useful carrier that can provide solubility for any pH modifier and/or provide good food soil pick up and/or have sufficiently low viscosity upon heating and/or allows the solid cleaner to retain its shape at room temperature. In many embodiments, the carrier is generally recognized as safe (GRAS) for food contact. Carriers include, for example, water, glycerin, triethylene glycol, and diethylene glycol. In some embodiments, the carrier is present in the solid cleaner in range from 0 to 80 wt %, or from 1 to 60 wt %, or from 25 to 50 wt %. [0047] In some embodiments, the carrier includes glycerin or glycerol. In certain embodiments, glycerin or glycerol can also act as a solubilizer of soils to be cleaned from the heated surfaces. When present, glycerin can make up from 1 to 80 wt %, or from 1 to 50 wt %, or from 5 to 40 wt %, or from 10 to 30 wt %. In some embodiments, the carrier includes water. When present, water can make up from 1 to 80 wt %, or from 1 to 50 wt %, or from 5 to 40 wt %, or from 10 to 30 wt %. In further embodiments, the carrier includes water and glycerin. When present, water and glycerin can make up from 1 to 80 wt %, or from 1 to 50 wt %, or from 5 to 40 wt %, or from 10 to 30 wt %. [0048] Thickeners can be optionally included in the solid cleaner, as desired. In many embodiments, thickeners can replace a portion of the solidifying agent, as desired. Thickeners can include, for example, xanthan gum, guar gum, polyols, alginic acid, sodium alginate, propylene glycol, methyl cellulose, polymer gels, clay, gelatin/clay mixtures, gelatin/oxide nanocomposite gels, smectite clay, montmorillonite clay, fillers e.g. CaCO 3 and mixtures of therein. If present, thickeners can make up from 0.1 to 25 wt %, or from 0.5 to 10 wt %. [0049] Abrasive material can be optionally included in the solid cleaner, as desired. In many embodiments, the abrasive materials incorporated into the solid cleaning composition can assist in the mechanical scrubbing action and can be used alone or in addition to an abrasive pad described herein. Abrasive materials include, for example, inorganic abrasive particles, organic based particles, sol gel particles or combinations thereof. Further examples of suitable abrasive particles are described in WO 97/49326. [0050] Additives can be optionally included in the solid cleaner, as desired. Additives can include, for example, builders, corrosion inhibitors (e.g., sodium benzoate), sequestering agents (EDTA), dyes, preservatives, and fragrances. In many embodiments, the additives are generally recognized as safe (GRAS) for food contact or approved for use as a direct food additive. [0051] In some embodiments, a non-woven substrate can be combined with the solid cleaners disclosed herein. Non-woven substrates are suited for scouring heated surfaces and can assist in physical removal of food soils at least partially removed or softened by the solid cleaners disclosed herein. In many embodiments, non-woven substrates include non-woven webs of fibers. [0052] In general, non-woven webs of fibers may be made of an air-laid, carded, stitch-bonded, thermobonded and/or resin-bonded construction of fibers, all as known by those skilled in the art. Fibers suitable for use in non-woven substrate materials include natural and synthetic fibers, and mixtures thereof. Synthetic fibers are preferred including those made of polyester (e.g., polyethylene terephthalate), nylon (e.g.; hexamethylene adipamide, polycaprolactam), polypropylene, acrylic (formed from a polymer of acrylonitrile), rayon, cellulose acetate, and so forth. Suitable natural fibers include those of cotton, wool, jute, and hemp. The fiber material can be a homogenous fiber or a composite fiber, such as bicomponent fiber (e.g., a co-spun sheath-core fiber). Non-woven substrate materials may also include different fibers in different portions. In some thermobonded non-woven substrate embodiments, the substrate includes melt bondable fibers where the fibers are bonded to one another by melted portions of the fibers. [0053] In some embodiments, the non-woven substrate material is an open, low density, three-dimensional, non-woven web of fibers, the fibers bonded to one another at points of mutual contact, referred to in the following as a “lofty, nonwoven web material”. In some embodiments, the fibers are thermo-bonded and/or resin-bonded (i.e. with a hardened resin, e.g. a prebond resin) to one another at points of mutual contact. In other embodiments, the fibers are resin-bonded to one another at points of mutual contact. Because the fibers of the web are bonded together at points of mutual contact, e.g. where they intersect and contact one another, a three-dimensional web structure of fibers is formed. The many interstices between adjacent fibers remain substantially unfilled, for example by resin, and thus an open web structure of low density having a network of many relatively large intercommunicated voids is provided. The term “open, low density” non-woven web of fibers is understood to refer to a non-woven web of fibers that exhibits a void volume (i.e. percentage of total volume of voids to total volume occupied by the non-woven web structure) of at least 75%, or at least 80%, or at least 85%, or in the range of from 85% to at least 95%. Such a lofty, non-woven web material is described in U.S. Pat. No. 2,958,593, which is incorporated by reference herein. [0054] Another example of a lofty, non-woven web material is described by U.S. Pat. Nos. 2,958,593, and 4,227,350, which are incorporated by reference herein. These patents disclose a lofty, non-woven web formed from a continuous extrusion of nylon coil material having a diameter in a range from 100 micrometers to 3 mm. Inorganic and/or organic abrasive materials can be optionally included on these non-woven webs. [0055] In some resin-bonded, lofty non-woven web material embodiments, the resin includes a coatable resinous adhesive such as a thermosetting water based phenolic resin, for example. Polyurethane resins may also be employed as well as other resins. Those skilled in the art will appreciate that the selection and amount of resin actually applied can depend on any of a variety of factors including, for example, fiber weight, fiber density, fiber type as well as the contemplated end use. Suitable synthetic fibers for production of such a web include those capable of withstanding the temperatures at which selected resins or adhesive binders are cured without deterioration. [0056] In some lofty, non-woven web material embodiments, suitable fibers are between 20 and 110 mm, or between 40 and 65 mm, in length and have a fineness or linear density ranging from 1.5 to 500 denier, or from 1.5 to 100 denier. Fibers of mixed denier can also be used, as desired. In one embodiment, the non-woven substrate includes polyester or nylon fibers having linear densities within the range from 5 to 65 denier. [0057] Lofty, non-woven web materials may be readily formed, e.g. air laid, for example, on a “Rando Webber” machine (commercially available from Rando Machine Company, New York) or may be formed by other conventional processes such as by carding or by continuous extrusion. Useful lofty, non-woven substrate materials have a fiber weight per unit area of at least 25 g/m 2 , or at least 50 g/m 2 , or between 50 and 1000 g/m 2 , or between 75 and 500 g/m 2 . Lesser amounts of fiber within the lofty, non-woven substrate materials will provide webs, which may be suitable in some applications. [0058] The foregoing fiber weights will provide a useful non-woven substrate having a thickness from 5 to 200 mm, or between 6 to 75 mm, or between 10 and 30 mm. For phenolic prebond resins applied to a lofty, non-woven substrate having a fiber weight within the above ranges, the prebond resin is applied to the web or substrate in a relatively light coating, providing a dry add-on weight within the broad range from 50 to 500 g/m 2 . [0059] The foregoing lofty, non-woven substrate materials are effective for most scouring applications. For more intensive scouring applications, the lofty, non-woven substrate materials may be provided with abrasive particles dispersed and adhered there within. The abrasive particles can be adhered to the surfaces of the fibers in the lofty, non-woven substrate material. In many embodiments, the abrasive particles may include inorganic abrasive particles, organic based particles, sol gel particles or combinations thereof, all as known in the art. Examples of suitable abrasive particles as well as methods and binders for adhering abrasive particles onto the surfaces of the fibers are for example described in WO 97/49326. [0060] In some embodiments, abrasive particles are adhered to the fibers of the non-woven substrate by a hardened organic resin binder such as, for example, a heat cured product of a thermosetting coatable resinous adhesive applied to the fibers of the non-woven substrate as a “binder precursor”. As used herein, “binder precursor” refers to a coatable resinous adhesive material applied to the fibers of the non-woven substrate to secure abrasive particles thereto. “Binder” refers to the layer of hardened resin over the fibers of the nonwoven web formed by hardening the binder precursor. In some embodiments, the organic resins suitable for use as a binder precursor in the non-woven substrate are formed from an organic binder precursor in a flowable state. During the manufacture of the non-woven substrate, the binder precursor can be converted to a hardened binder or make coat. In some embodiments, the binder is in a solid, non-flowable state. In some embodiments, the binder is formed from a thermoplastic material. In other embodiments, the binder is formed from a material that is capable of being cross-linked. In some embodiments, a mixture of a thermoplastic binder and a cross-linked binder is also useful. [0061] During the process to make the web or substrate, the binder precursor can be mixed with the foregoing abrasive particles to form an adhesive/abrasive slurry that may be applied to the fibers of the non-woven by any of a variety of known methods such as roll coating, knife coating, spray coating, and the like. The thus applied binder precursor is then exposed to the appropriate conditions to solidify the binder. For cross-linkable binder precursors, the binder precursor can be exposed to the appropriate energy source to initiate polymerization or curing and to form the hardened binder. [0062] In some embodiments, the organic binder precursor is an organic material that is capable of being cross-linked. The binder precursors can be either a condensation curable resin or an addition polymerizable resin, among others. The addition polymerizable resins can be ethylenically unsaturated monomers and/or oligomers. Examples of useable cross-linkable materials include phenolic resins, bis-maleimide binders, vinyl ether resins, aminoplast resins having pendant alpha,beta-unsaturated carbonyl groups, urethane resins, epoxy resins, acrylate resins, arylated isocyanurate resins, urea-formaldehyde resins, melamine formaldehyde resins, phenyl formaldehyde, styrene butadiene resins, isocyanurate resins, acrylated urethane resins, acrylated epoxy resins, or mixtures thereof. The binder precursor suitable for use is a coatable, hardenable adhesive binder and may comprise one or more thermoplastic or, thermosetting resinous adhesives. Resinous adhesives suitable for use in the present invention include phenolic resins, aminoplast resins having pendant alpha,beta-unsaturated carbonyl groups, urethane resins, epoxy resins, ethylenically unsaturated resins, acrylated isocyanurate resins, urea-formaldehyde resins, isocyanurate resins, acrylated urethane resins, acrylated epoxy resins, bismaleimide resins, fluorine-modified epoxy resins, and combinations thereof. Examples of these resins can be found in WO 97/49326. Catalysts and/or curing agents may be added to the binder precursor to initiate and/or accelerate the polymerization process. In many embodiments the substrate can withstand temperatures up to at least 200 degrees Celsius, (e.g., food preparation operating temperature.) [0063] Commercially available non-woven substrate or web materials are available under the trade designation “Scotch-Brite™ General Purpose Scour Pad No. 96,” “Scotch-Brite™ Heavy Duty Griddle Cleaner No. 82 (non-woven glass cloth),” “Scotch-Brite™ All Purpose Scour Pad No. 9488R,” “Scotch-Brite™ Heavy Duty Scour Pad No. 86,” all available from 3M Co. In other embodiments, the substrate is a Scotch-Brite™ Griddle Screen No. 68, a Scotch-Brite™ Griddle Screen No. 200, steel-wool, pumice block, foamed glass bricks, and the like. EXAMPLES [0064] All chemicals were used as commercially available. [0000] Table of Abbreviations Abbreviation Description Quick Clean Scotch-Brite ™ Quick Clean Griddle Liquid, No. 700, 3M Co., St. Paul, MN FAME Fatty Acid Mono Ester (Lauricidin ™), Med-Chem. Laboratories, Galena, IL PEG Poly(ethylene glycol) (1000 Da, 4600 Da, or 8000 Da), Aldrich, Milwaukee, WI. Potassium Carbonate Ashta Chemicals, Ashtabula, OH. K 2 CO 3 (anhydrous) Sodium Carbonate J. T. Baker, Phillipsburg, NJ. Na 2 CO 3 (monohydrate) Stock Solution #1 10 g Potassium Carbonate/4 g Sodium Carbonate/20 g DI Water Stock Solution #2 12 g Potassium Carbonate/6 g Sodium Carbonate/20 g DI Water Stock Solution #3 10 g Potassium Carbonate/4 g Sodium Carbonate/15 g DI Water Stock Solution #4 10 g Potassium Carbonate/4 g Sodium Carbonate/14 g DI Water Glycerin Merck KGaA, Darmstadt Germany TONE Polyol 210 Melting Point Range: 35° to 45° C., Dow/Union Carbide, Midland, MI TONE Polyol 230 Melting Point Range: 40° to 50° C., Dow/Union Carbide, Midland, MI TONE Polyol 240 Melting Point Range: 45° to 55° C., Dow/Union Carbide, Midland, MI TONE Polyol 260 Melting Point Range: 50° to 60° C., Dow/Union Carbide, Midland, MI #46 Pad Scotch-Brite ™ Griddle Polishing Pad No. 46, 3M Co., St. Paul, MN #9488R Pad Scotch-Brite ™ All Purpose Scouring Pad No. 9488R, 3M Co., St. Paul, MN SPAN 40 Sorbitan Monopalmitate Surfactant, Aldrich, Milwaukee, WI SPAN 65 Sorbitan Tristearate Surfactant, Imperial Chemical Industries (ICI), London, UK Brij 35 Dodecylpoly(ethylene glycol) ether surfactant, Uniquema (ICI), London, UK Pluracare L44 NF Block copolymer of poly(ethylene glycol) and poly(propylene glycol), BASF, Lundwigshafen, DE BioSoft D-40 Sodium Dodecylbenzene Sulphonate Surfactant, Stepan Company, Northfield, IL EDTA Ethylene Diamine Tetra Acetate - Sequesterant Eastman Kodak Co., Kingsport, TN Xanthan Gum R. T. Vanderbilt Company, Inc. Norwalk, CT. Candelilla wax Strahl & Pitsch, Inc., West Babylon, CT. Sodium Metasilicate J. T. Baker, Phillipsburg, NJ. Sodium Bicarbonate Mallinckrodt BaKER, Inc., Paris, KY Melamine Particle 40/100 mesh. Maxi-Blast, Inc., South Bend, IN. formaldehyde particles Pumice 0 Charles B. Chrystal Co., Inc. New York, NY Pumice FF Charles B. Chrystal Co., Inc. New York, NY Emulsifying wax NF Strahl & Pitsch, Inc., West Babylon, CT. Cetyl Alcohol TCI Mark Stearyl Alcohol Alfol 18 - Sasol North America Inc., Weslake, Louisiana. Test Methods for Cleaning the Griddle Burnt Oil Test Method [0000] 1. Turn all three burners on the flat griddle (Star Mftg. Model 536-76A. Smithville Tenn.) to 450° F. (232° C.). 2. Measure about 40 mL of commercially available soybean oil (e.g., Crisco) and pour on the griddle. 3. Spread out oil with a 3M Green Scotch-Brite™ General Purpose Scour Pad No. 96 until even over entire surface of griddle. 4. Let griddle heat oil for 45 minutes. Oil should be dark brown and of fairly uniform color across the entire griddle. 5. Decrease the temperature of the griddle to 300-350° F. (150-175° C.). 6. Measure the temperature of the griddle with the IR thermometer (Dickson, Chicago, Ill.) and record it. It should be between 300-350° F. (150-175° C.). 7. Apply test cleaning composition on desired amount of griddle. 100 grams of test cleaning composition for the entire griddle. 8. Apply test cleaner over griddle surface with Scotch-Brite™ Griddle Polishing Pad No. 46 on pad holder and record the amount of time for the entire product to melt. 9. Turn off burner under section of griddle you are testing. 10. Immediately begin scrubbing using #46 pad and record amount of time necessary for acceptable level of cleanliness. 11. Scrape griddle surface with squeegee to move melted wax into grease trap. 12. Repeat cleaning over other surfaces of griddle with other test cleaners. 13. Using a wet paper towel on the pad holder, rinse surface and edges of griddle. 14. Apply a small amount of oil to surface of griddle and spread with Scotch-Brite™ General Purpose Scour Pad No. 96 to season the surface. 15. Wipe up any excess oil with a paper towel Ground Beef Test Method [0000] 1. Turn all three burners to 325° F. (160° C.). 2. Weigh 2.5 lbs (1.1 Kg) of ground beef for the entire griddle 3. Cook the beef until dark brown, moving the ground beef around the griddle to make it evenly distributed. 4. Remove the beef from the griddle with the flat cooking utensil taking off as much beef as possible. 5. Leave the food soil cooking for an extra 60 minutes 6. Measure the temperature of the griddle and record it. It should be between 300-350° F. (150-175° C.). 7. Apply test cleaner over desired amount of griddle. 100 g to 120 g of cleaning composition for the entire griddle. 8. Spread test cleaner over griddle surface with an appropriate pad (either 3M #46 Griddle Polishing Pad or 3M #9488R All Purpose Pad) on pad holder and record the amount of time for the entire product to melt. 9. Turn off burner under section of griddle you are testing. 10. Immediately begin scrubbing using the No. 46 pad and record amount of time necessary for acceptable level of cleanliness. 11. Scrape griddle surface with squeegee. 12. Repeat cleaning over the entire surfaces of griddle with other test cleaners. 13. Using a wet paper towel on the pad holder, rinse surface and edges of griddle. 14. Wash out drip tray of any remaining food soil. 15. Apply a small amount of oil to surface of griddle and spread with Scotch-Brite™ General Purpose Scour Pad No. 96 to season to surface. 16. Wipe up any excess oil with a paper towel. [0096] Preparation of the Cleaning Compositions [0097] Stock solutions were made by dissolving the salts indicated below in de-ionized water at low heat. The solution was stirred until no more solid salts were present. [0098] The stock solutions and glycerin (Procter & Gamble, Cincinnati, Ohio) were added to a beaker and placed on a hot plate/stirrer. The solution was heated to about 80° C. while gently mixing. The solidifying agent (wax or polyol) was added to the stock solution/glycerin mix and heated while stirring until the solidifying agent was completely melted. The formulation was taken off the heat once it was well mixed and homogenous. [0099] Tablets and impregnated pads were made by either pouring into the molds to form tablets or pads. Tablets were made by allowing the melted formulations to cool down to room temperature in an aluminum mold of 2″×2″×1″ (5 cm×5 cm×2.5 cm) (W×L×H). Tablets of 60 g each were made with this mold. Impregnated pads (#46) were also made by pouring the melted formulation on a mold of 4″×5″×1″ (10 cm×13 cm×2.5 cm) (W×L×H) at about 80° C., allowing it to cool down to about 60° C. and then placing the pad onto the mold and applying a little pressure to force the pad into the solidified cleaner. The pads were allowed to cool to room temperature. [0100] Formulations were also made of the following waxes: Rice bran wax (Koster Keunen, Inc., Watertown, Conn., USA) Lemon peel Wax (Koster Keunen, Inc., Watertown, Conn., USA) Soy wax flakes (Koster Keunen, Inc., Watertown, Conn., USA) Deodorized orange peel wax (Koster Keunen, Inc., Watertown, Conn., USA) Beeswax (Strahl & Pitsch, Inc., West Babylon, N.J., USA) Candelilla wax (Strahl & Pitsch, Inc., West Babylon, N.J., USA) Carnauba wax (Strahl & Pitsch, Inc., West Babylon, N.J., USA) [0108] Formulation 1 [0109] A solid cleaner was made by combining 34 g of stock solution #1 with 22 g of glycerin and 44 g of beeswax. [0110] Formulation 2 [0111] A solid cleaner was made by combining 34 g of stock solution #1 with 22 g of glycerin and 44 g of carnauba wax. [0112] Formulation 3 [0113] A solid cleaner was made by combining 34 g of stock solution #1 with 22 g of glycerin and 44 g of candelilla wax. [0114] Formulation 4 [0115] A solid cleaner was made by combining 34 g of stock solution #1 with 33 g of glycerin and 33 g of beeswax. [0116] Formulation 5 [0117] A solid cleaner was made by combining 34 g of stock solution #1 with 33 g of glycerin and 33 g of carnauba wax. [0118] Formulation 6 [0119] A solid cleaner was made by combining 34 g of stock solution #1 with 40 g of glycerin and 26 g of carnauba wax. Formulation 7 [0120] A solid cleaner was made by combining 34 g of stock solution #1 with 40 g of glycerin and 26 g of candelilla wax. [0121] Formulation 8 [0122] A solid cleaner was made by combining 34 g of stock solution #2 with 40 g of glycerin and 26 g of candelilla wax. [0123] Formulation 9 [0124] A solid cleaner was made by combining 34 g of stock solution #2 with 40 g of glycerin and 26 g of candelilla wax impregnated into a pad. [0125] Formulation 10 [0126] A solid cleaner was made by combining 34 g of stock solution #2 with 40 g of glycerin and 26 g of beeswax impregnated into a pad. [0127] Formulation 11 [0128] A solid cleaner was made by combining 34 g of stock solution #2 with 40 g of glycerin and 26 g of carnauba wax impregnated into a pad. [0129] Formulation 12 [0130] A solid cleaner was made by combining 34 g of stock solution #2 with 40 g of glycerin and 26 g of lemon peel wax. [0131] Formulation 13 [0132] A solid cleaner was made by combining 24 g of stock solution #2 with 40 g of glycerin and 26 g of carnauba wax and 10 g of sodium bicarbonate. [0133] Formulation 14 [0134] A solid cleaner was made by combining 24 g of stock solution #2 with 40 g of glycerin and 26 g of carnauba wax and 10 g of sodium metasilicate. [0135] Formulation 15 [0136] A solid cleaner was made by combining 34 g of stock solution #2 with 40 g of glycerin and 26 g of rice wax. [0137] Formulation 16 [0138] A solid cleaner was made by combining 34 g of stock solution #2 with 40 g of glycerin and 26 g of orange peel wax. Results [0139] Experimental samples were compared against Scotch-Brite™ Quick Clean Griddle Liquid No. 700 (Quick Clean or 700) (3M Company, St. Paul, Minn.) and rated for melting time (in seconds), and cleaning performance. A visual rating was given for cleaning performance. The rating scale went from 1 to 5, with 5 being no food residue left on the heated surface. The temperature of the griddle was recorded with an IR thermometer. [0140] A comparison of the performance of the different experimental formulations against Quick Clean is shown in the table below. [0000] Griddle Cleaner Evaluation Griddle Melting Temperature Time Cleaning Example Formulation Soil (° F.) (sec) Performance 1 1 Oil — — 3 2 2 Oil — — 3 3 3 Oil — — 3 4 4 Oil 330 (165° C.) 38 3 5 5 Oil 325 (160° C.) 45 3 6 6 Oil 300 (150° C.) 42 3 7 Quick Clean Oil 330 (165° C.) N/A 5 8 7 Oil 330 (165° C.) 40 3 9 8 Oil 325 (160° C.) 42 5 10 9 Oil 330 (165° C.) — 5 11 9 Oil 325 (160° C.) 110 5 12 10 Oil 335 (168° C.) 40 5 13 11 Oil 325 (160° C.) 30 3 14 8 Beef 350 (175° C.) 85 5 15 8 Beef 350 (175° C.) 120 5 16 8 Beef 360 (182° C.) 19 5 17 8 Beef 360 (182° C.) 67 5 18 Quick Clean Beef 340 (171° C.) N/A 5 19 11 Oil 350 (175° C.) 45 5 20 12 Oil 340 (171° C.) 54 5 21 15 Oil 330 (165° C.) 38 5 22 16 Oil 325 (160° C.) 32 3 Further Prepared and Tested Samples: [0141] The following formulations were made up using Quick Clean, FAME, PEG 1000, 4600 and 8000 as well as Stock Solutions #1 and #3 (defined in the Table of Abbreviations above). [0000] Compositions in % wt PEG Stock Solution Example # FAME 1000 4600 8000 #1 #3 Quick Clean (1) — — — — — — 23 16 — 50 — — 34 24 16 — — 50 — 34 25 36 30 — — — 34 26 36 — — — — 34 27 36 — 30 60 — 34 28 50 16 — — — 34 29 50 — 16 — 34 — 30 50 — 16 — — 34 31 50 — — 16 — 34 [0142] The following formulations were made up using Glycerin, TONE Polyols (210, 230, 240 and 260), Stock Solution #3 (defined in the Table of Abbreviations above). In addition, Example #42 and #43 were loaded into a Scotch-Brite™ Griddle Polishing Pad No. 46. [0000] Composition in % wt Stock Example Difunctional TONE Polyol solution Loaded # Glycerin 210 230 240 260 #1 #3 Pad 32 13 69 — — — — 18 NO 33 13 — 69 — — — 18 NO 34 13 — — 69 — — 18 NO 35 13 — — — 69 — 18 NO 36 13 69 — — — — 18 YES 37 13 — — — 69 — 18 YES [0143] The following formulations were made up using Glycerin, TONE Polyols (210 and 260), SPAN 40, SPAN 65, Quick Clean and Stock Solutions #3 and #4 (defined in the Table of Abbreviations above). [0000] Composition in % wt Difunctional TONE Surfactant Polyol SPAN SPAN Quick Stock Solution Example # Glycerin 210 260 40 65 Clean #3 #4 38 13 — 61 10 — — 16 — 39 13 — 61 — 10 — 16 — 40 — — 77 — — 23 — — 41 13 41 33 — — — 13 — 42 13 67 — — — — — 20 43 13 — 68 — — — — 19 [0144] The following formulations were made up using Glycerin, TONE Polyols (210 and 260), SPAN 40, Brij 35, Pluracare L44 NF, BioSoft D-40, PEG 1000, and Stock Solution #3 (defined in the Table of Abbreviations above). [0000] Composition in % wt Difunctional Surfactants/Detergents Stock TONE Polyol Span Brij Pluracare BioSoft PEG Sol. Example # Glycerin 210 260 40 35 L44 NF D-40 1000 #3 44 14 68 — — 0.05 — — — 18 45 14 68 — — — — 0.2 — 18 46 13 69 — — — 0.05 — — 17 47 14 58 — — — — — 10 16 48 11 — 66 — — — — 8 14 49 14 67 — 1 — — — — 18 50 14 — 67 1 — — — — 18 [0145] The following formulations were made up using Quick Clean, Glycerin, TONE Polyols (210 and 260), SPAN 40, EDTA, and Stock Solution #2 (defined in the Table of Abbreviations above). [0000] Composition in % wt Difunctional Stock Example TONE Polyol Surfactant Sequester Sol. # Glycerin 210 260 SPAN 40 EDTA #3 51 14 — 66 — 3 17 52 14 67 — 0.05 3 17 53 13 71 — 0.05 1 15 [0146] The following griddle cleaner formulations were made using Stock Solution #2, Glycerin, Candelilla Wax, and Xanthan Gum. The stock solution and glycerin were added to a beaker and placed on a hot plate/stirrer. The solution was heated to about 100° C. while gently mixing. The wax was added to the stock solution/glycerin mix and left in the heat while stirring until the wax was completely melted. Xanthan gum was added to the formulations at 100° C. after the wax was melted. The formulation was taken off the heat once it was well mixed and homogeneous. [0147] Tablets and impregnated pads were made by either pouring into the molds to form tablets or pads. Tablets were made by allowing the melted formulation to cool down to room temperature in an aluminum mold of 2″×2″×1″ (5 cm×5 cm×2.5 cm) (W×L×H). Tablets of 50 g each were made with this mold. Impregnated pads (#46) were also made by pouring the melted formulation on a mold of 4″×5.5″×1″ (10 cm×14 cm×2.5 cm) (W×L×H) at about 80° C., allowing it to cool down to about 60° C. and then placing the pad and applying a little pressure. Pads of 100 g each were allowed to cool to room temperature. [0000] Stock Candelilla Xanthan Solution Glycerin Wax Gum Example # #2 (g) (g) (g) (g) 54 42.7 41.0 16.3 0.0 55 42.2 40.4 16.1 1.2 56 40.2 38.5 15.4 5.9 57 39.3 37.6 15.0 8.1 58 50.0 29.4 19.1 1.5 59 47.2 27.8 18.1 6.9 Formulation 9 34.0 40.0 26.0 0.0 [0148] Performance of these examples were compared to the control sample Formulation 9 (solid cleaner with no xanthan gum). Formulations were rated for cleaning performance. A visual rating was given for each of these qualitative attributes listed above. The rating scale went from 1 to 5, with 5 being best. [0000] Stock Solution #2 Glycerin Candelilla Xanthan ratio Melting time Cleaning Example # (g) (g) Wax (g) Gum (g) Gly/Wax (sec) performance 54 42.7 41.0 16.3 0.0 2.5 45 5 55 42.2 40.4 16.1 1.2 2.5 50 5 56 40.2 38.5 15.4 5.9 2.5 40 5 57 39.3 37.6 15.0 8.1 2.5 40 1 58 50.0 29.4 19.1 1.5 1.5 38 4 59 47.2 27.8 18.1 6.9 1.5 36 1 Formulation 34.0 40.0 26.0 0.0 1.5 45 5 9 [0149] Results appear to indicate that formulations containing xanthan gum up to 6% were solid even when the amount of candelilla wax was significantly reduced from 26 g to 15-16 g. Examples 55 and 56 appear to show performance comparable to that of the control sample Formulation 9 (formulation with no thickener and higher wax content). [0150] A variety of abrasive materials were added to Formulation 9 to form the Examples listed in the table below. The examples including abrasive materials were loaded onto the non-abrasive #9488R pad, while the Formulation 9 and the quick clean example was loaded onto an abrasive #46 pad. Tablets and impregnated pads were made by either pouring into the molds to form tablets or pads. Tablets were made by allowing the melted formulation to cool down to room temperature in an aluminum mold of 2″×2″×1″ (5 cm×5 cm×2.5 cm) (W×L×H). Tablets of 50 g each were made with this mold. Impregnated pads were also made by pouring the melted formulation on a mold of 4″×5.5″×1″ (10 cm×14 cm×2.5 cm) (W×L×H) at about 80° C., allowing it to cool down to about 60° C. and then placing the pad and applying a little pressure. Pads of 100 g each were allowed to cool to room temperature. [0151] Performance of these examples were compared to the control sample Formulation 9 (solid cleaner with no abrasive) and to Quick Clean. Formulations were rated for cleaning performance. A visual rating was given for each of these qualitative attributes listed above. The rating scale went from 1 to 5, with 5 being best. [0000] Grams of Cleaning Abrasive/ perfor- Example # Abrasive 100 g of Wax Soil mance 60 Sodium Bicarbonate 10 Oil 1 61 Sodium Bicarbonate 20 Oil 5 62 Sodium Metasilicate 10 Oil 1 63 Sodium Metasilicate 20 Oil 1 64 Pumice 0 10 Oil 3 65 Pumice 0 20 Oil 4 66 Pumice 0 30 Oil 1 67 Pumice 0 50 Oil 1 68 Pumice FF 10 Oil 3 69 Pumice FF 20 Oil 4 70 Pumice 0 10 Beef 5 71 Pumice FF 10 Beef 5 72 Melamine Resin 10 Oil 5 73 Melamine Resin 20 Oil 5 74 Melamine Resin 30 Oil 5 Formulation 9 — — Oil 5 Quick Clean — — Oil 5 Formulation 9 — — Beef 5 Quick Clean — — Beef 5 [0152] These results appear to indicate that the performance of abrasive containing formulations was the same or better than the Quick Clean and control sample Formulation 9. [0153] Emulsifying Wax NF was added to Formulation 9 to form the Examples listed in the table below. Tablets and impregnated pads were made by either pouring into the molds to form tablets or pads. Tablets were made by allowing the melted formulation to cool down to room temperature in an aluminum mold of 2″×2″×1″ (5 cm×5 cm×2.5 cm) (W×L×H). Tablets of 50 g each were made with this mold. Impregnated pads (#46) were also made by pouring the melted formulation on a mold of 4″×5.5″×1″ (10 cm×14 cm×2.5 cm) (W×L×H) at about 80° C., allowing it to cool down to about 60° C. and then placing the pad and applying a little pressure. Pads of 100 g each were allowed to cool to room temperature. [0154] Performance of these examples were compared to the control sample Formulation 9 (solid cleaner with no emulsifying wax). Formulations were rated for cleaning performance. A visual rating was given for each of these qualitative attributes listed above. The rating scale went from 1 to 5, with 5 being best. [0000] Stock Solution #2 Glycerin Candelilla Emulsifying ratio Melting Cleaning Example # (g) (g) Wax (g) Wax NF (g) Cand/Emul time (sec) performance 75 34 40 13 13 1:1 25 5 76 34 40 9 17 1:2 30 5 77 34 40 17 9 2:1 30 5 78 34 40 20 6 3:1 35 5 Formulation 34 40 26 0 0 45 5 9 79 34 30 13 13 1:1 30 5 80 34 25 13 13 1:1 25 5 81 34 20 13 13 1:1 25 5 [0155] These results appear to indicate that formulations that contain Emulsifying Wax NF melt faster than the control sample formulation 9. In addition, formulations that contain Emulsifying Wax NF were reported to have less “drag” when applied to the heated surface than the control sample formulation 9. [0156] The following formulations were made up using stock solution #2, glycerin, wax and an emulsifying wax (cetyl and/or stearyl alcohol). [0000] Stock Solution #2 Glycerin Candelilla Carnauba Cetyl Stearyl Melting Cleaning Example # (g) (g) Wax (g) Wax (g) Alcohol (g) Alcohol (g) time (sec) performance 82 34 40 13 0 0 13 38 5 83 34 40 13 0 13 0 35 5 84 34 40 13 0 6.5 6.5 38 5 85 34 40 0 13 0 13 48 5 86 34 30 0 13 0 13 33 5 [0157] All references and publications cited herein are expressly incorporated herein by reference in their entirety into this disclosure. Illustrative embodiments of this disclosure are discussed and reference has been made to possible variations within the scope of this disclosure. These and other variations and modifications in the disclosure will be apparent to those skilled in the art without departing from the scope of this disclosure, and it should be understood that this disclosure is not limited to the illustrative embodiments set forth herein. Accordingly, the disclosure is to be limited only by the claims provided below. [0158] Referring to FIGS. 3 a - c and 4 a - c , alternative embodiments of the pad 20 shown in FIGS. 1 and 2 are shown. The pad 40 includes a stepped cross sectional profile that is different than the cross sectional profile of pad 20 . The pad 40 includes valleys 42 separated by peaks 44 , however, the pad 40 includes angled edge surfaces 46 and 48 that slope away from the flat surface 50 of the peaks 44 down towards the flat surface 52 of the valleys 42 . Pad 40 may be preferred over the stepped profile of pad 20 for some griddle surface configurations. For example, in contrast to the griddle 34 shown in FIG. 1 , which has straight vertical edge surfaces 54 and 56 (shown in FIG. 1 ), as shown in FIG. 3 b other griddle configurations include sloped surfaces 60 and 62 that connect the high portions 64 of the griddle surface to the low portions 66 of the griddle surface. In addition, other griddle configurations may include curved top surfaces 70 and curved bottom surfaces 72 that are connected by curved side surfaces 74 and 76 . For such griddle configurations, the pad 40 may be preferred. [0159] Still referring to FIGS. 3 a - c and 4 a - c , the stepped profile of pad 40 may also be preferred in situations where the pad 40 is expected to be used on griddle surfaces having unknown or variable grooved spacing GGS. The flat surface 50 of the pad 40 can be configured such that it will fit between grooves even on griddles having grooves that are relatively close together. In one embodiment the width W 1 of the flat surface 50 is set to fit in the grooves of griddles having the smallest griddle spacing GGS, and the distance D 2 from the center of one valley to the next is set to accommodate the raised portion of griddles having the largest griddle spacing GGS. In such embodiments the distance between the center of two adjacent valleys D 2 may be greater than twice the width W 1 of the flat surface 50 . [0160] The pad 40 in the depicted embodiment is geometrically configured such that a single model can work well to clean a number of different commercially available griddles having different surface configurations. While in use the pad 40 can be moved back and forth along the griddle grooves in the X-direction while biased against the right side 78 of the peaks 80 in the positive Y-direction (shown in FIG. 4 b ) to clean the first portion 82 of the griddle surface. Next, the pad 40 can be moved back and forth along the grooves in the X-direction while biased against the left side 84 of the peaks 80 in the negative Y-direction (shown in FIG. 4 c ) to clean the second portion 86 of the griddle surface. [0161] Referring to FIGS. 5 and 6 , another embodiment of the pad according to the invention is shown. The pad 90 includes a plurality of separate pad sections 92 , 94 , and 96 that are held together by metal wire loops 98 and 100 . The metal wires loops 98 and 100 extend through a center portion 102 of each of the separate pad sections 92 , 94 , and 96 , thereby holding the center portions 102 of each pad section 92 , 94 and 96 together. The upper end portions 104 and lower end portions 106 of each pad section 92 , 94 , and 96 are free to deflect a small distance in the X-direction relative to each other. The capability of the cleaning end 106 or 104 to deflect can enable the pad 90 to be more compatible with griddles having different groove spacing GGS and different surface profiles. As pressure is applied to the pad 90 the pad deforms such that the pad 90 matches the surface profile of whatever griddle surface configuration it is position over. It should be appreciated that many other suitable materials may be used in place of metal loop 98 and 100 to hold the pad sections 92 , 94 , and 96 together. For example, in an alternative embodiment nylon straps may be used in place of the metal wires 98 and 100 . It should also be appreciated that any number of strap configurations can be used to band the pad together. In other words, the device that holds the pads 92 , 94 , and 96 together need not be looped. For example, in other embodiments the pads 92 , 94 , and 96 may be stapled together, heat staked together, ultrasonically bonded, or glued together. In should also be appreciated that though only three pad sections 92 , 94 , and 96 are shown, any number of pad sections may be used to form the complete pad 90 . [0162] Referring to FIG. 7 , another embodiment of the pad according to the invention is shown. The pad 110 includes preformed creases 112 , 114 116 and 118 that enable the pad 110 to better fit the step profile bottom surface of the shoe 120 . The bottom of the shoe 120 can include any type of step profile desired. In the depicted embodiment the bottom surface 122 of the shoe 120 includes a plurality of hooks 124 that engage and secure the pad 110 thereto. It should be appreciated that though in the depicted embodiment the bottom surface 122 includes hooks 124 all across the bottom surface 122 , in alternative embodiments means other than hooks 124 may be used in attaching the pad 110 to the shoe 120 or possibly only particular areas of the bottom surface 122 may include hooks 124 . The creases 112 , 114 , 116 and 118 can be imparted onto the pad 110 by melting the pad along the creases to create a natural fold line in the pad 110 . Other means of creating the creases include scoring the pads along the fold lines. [0163] Referring to FIGS. 8 and 9 , another embodiment of the pad is shown. The pad 130 includes a number of pad members 132 - 148 that are stacked adjacent to each other and held together by a binding member 150 . The binding member 150 engages and secures the upper portions 152 of each pad 132 - 148 together to create a cleaning block. Relative to the upper portions 152 , the lower portions 154 of the pad members 132 - 148 are free to deflect. This deflection provides advantages in that the pad 130 can be used to clean a large variety of griddles having different surface geometries. When the pad 130 is pressed onto the griddle surface it conforms to fit the particular surface configuration of the griddle. In the embodiment shown each pad has a generally rectangular shape, but the block can be of any other shape as well. The upper portions 152 can be held together solely by the binding member 150 , or they can be glued or mechanically fastened together. For example, the metal wires 98 and 100 of the embodiment shown in FIGS. 5 and 6 can be used to hold the top portions 152 of the pads 132 - 148 together. [0164] Referring to FIGS. 8 and 9 , a method of assembling the pad 130 is shown. The method includes arranging pad members 132 - 148 adjacent each other and connecting the top portions of the pad members 132 - 148 together, then fitting the binding member 150 over the top portions 152 and around the pad members 132 - 148 . The binding member 150 includes an opening 154 that exposes portions of the upper edges 156 . The exposed portions of the upper edges engage the hooks 160 that extend from the foot portion 164 of the cleaning tool 162 . In the depicted embodiment the binding member 150 is a molded plastic part that is shaped like an open box frame with the center of the bottom of the box removed. In an alternative embodiment the binding member 150 could be constructed of a different material such as cardboard. In addition, many other ways to attach the pad 130 to the handle 162 are possible. [0165] Referring to FIGS. 10 and 11 the binding member 150 is shown in greater detail. The binding member includes four side surfaces 170 , 172 , 174 , 176 and top surface 178 . The top surface includes at least one opening 180 to allow the handles to engage the pad members (see FIG. 9 ). An alternative embodiment of the binding member 150 is shown in FIG. 12 . The binding member 182 includes a top surface 184 that has four openings 186 , 188 , 190 , and 192 instead of a single opening. In this embodiment the handle engages the pad members (see FIG. 9 ) through the four openings 186 , 188 , 190 , and 192 . The web portions 194 , 196 and 198 provide additional support for the pad members (see FIG. 9 ). [0166] The above specification provides a complete description of the manufacture and use of the composition of the invention. It should be understood that features from the depicted embodiments can be combined to form new embodiments not specifically depicted. Moreover, since many embodiments of the invention can be made without departing from the spirit and scope of the invention, the invention resides in the claims hereinafter appended.

A heated surface cleaning pad that can be attached to a bottom portion of a hand tool configured to be used to scrub a heated food preparation surface. According to an embodiment, the cleaning pad is configured to clean a hot griddle surface that has a number of parallel raised ridges. According to another embodiment, the cleaning pad includes a non-woven construction that is impregnated with abrasive particles and/or a liquid cleaner or a dissolvable solid cleaner. In addition, a method of cleaning a hot food preparation surface is provided. According to one embodiment, the method includes the steps of attaching a scrubbing pad to a cleaning tool, contacting the openings in the cleaning surface of the scrubbing pad with the raised ridges of a griddle surface, and scrubbing the griddle surface with the pad until the surface is clean. Moreover, a method of constructing a cleaning pad for cleaning a hot griddle surface is also provided. According to one embodiment, the method includes binding pad elements together such that a cleaning end of the cleaning pad is comprised of end portions of adjacently arranged scrubbing elements.

CROSS-REFERENCE TO RELATED APPLICATIONS [0001] This application is a divisional application of U.S. application Ser. No. 14/066,571, filed on Oct. 29, 2013, which is incorporated by reference in its entirety. TECHNICAL FIELD [0002] The present disclosure relates to chemical mechanical polishing and more specifically to monitoring of a conductive layer during chemical mechanical polishing. BACKGROUND [0003] An integrated circuit is typically formed on a substrate by the sequential deposition of conductive, semiconductive, or insulative layers on a silicon wafer. A variety of fabrication processes require planarization of a layer on the substrate. For example, one fabrication step involves depositing a filler layer over a non-planar surface and planarizing the filler layer. For certain applications, the filler layer is planarized until the top surface of a patterned layer is exposed. For example, a metal layer can be deposited on a patterned insulative layer to fill the trenches and holes in the insulative layer. After planarization, the remaining portions of the metal in the trenches and holes of the patterned layer form vias, plugs, and lines provide conductive paths between thin film circuits on the substrate. [0004] Chemical mechanical polishing (CMP) is one accepted method of planarization. This planarization method typically requires that the substrate be mounted on a carrier head. The exposed surface of the substrate is typically placed against a rotating polishing pad. The carrier head provides a controllable load on the substrate to push it against the polishing pad. Polishing slurry with abrasive particles is typically supplied to the surface of the polishing pad. [0005] One problem in CMP is determining whether the polishing process is complete, i.e., whether a substrate layer has been planarized to a desired flatness or thickness, or when a desired amount of material has been removed. Variations in the slurry composition, the polishing pad condition, the relative speed between the polishing pad and the substrate, the initial thickness of the substrate layer, and the load on the substrate can cause variations in the material removal rate. These variations cause variations in the time needed to reach the polishing endpoint. Therefore, determining the polishing endpoint merely as a function of polishing time can lead to non-uniformity within a wafer or from wafer to wafer. [0006] In some systems, a substrate is monitored in-situ during polishing, e.g., through the polishing pad. One monitoring technique is to induce an eddy current in the conductive layer and detect the change in the eddy current as the conductive layer is removed. SUMMARY [0007] In one aspect, a method of controlling polishing includes receiving a measurement of an initial thickness of a conductive film on a first substrate prior to polishing the first substrate from an in-line or stand-alone monitoring system, polishing one or more substrates in a polishing system, the one or more substrates including the first substrate, during polishing of the one or more substrates, monitoring the one or more substrates with an eddy current monitoring system to generate a first signal, determining a starting value of the first signal for a start of polishing of the first substrate, determining a gain based on the starting value and the measurement of the initial thickness, for at least a portion of the first signal collected during polishing of at least one substrate of the one or more substrates, calculating a second signal based on the first signal and the gain, and determining at least one of a polishing endpoint or an adjustment to a polishing parameter for the at least one of the substrate based on the second signal. [0008] Implementations may include one or more of the following features. [0009] Calculating the second signal may include multiplying by the gain. [0010] Calculating the second signal may include calculating V′=V*G+K wherein V′ is the second signal, V is the first signal, G is the gain and K is an offset. [0011] The at least one substrate of the one or more substrates may be the first substrate. [0012] The at least one substrate of the one or more substrates may be a second substrate polished subsequent to the first substrate. [0013] The polishing system may include a rotatable platen, and an eddy current sensor of the eddy current monitoring system is supported on the platen to sweep across the one or more substrates. [0014] The first signal may be generated from a portion of a signal generated when the eddy current sensor is not adjacent the substrate. [0015] A reference value may be determined from a portion of the signal generated when the eddy current sensor is not adjacent the substrate. [0016] An offset may be generated by adjusting the reference value to generate a desired value for zero thickness. [0017] Determining the gain may include determining a desired value from a calibration function relating thickness to signal strength and the measurement of the initial thickness. [0018] Determining the gain may include calculating a multiplier N according to [0000] N = ( D - K ) ( S - K ) [0019] where D is the desired value, S is the starting value, and K is a constant representing a value of the calibration function for zero thickness. [0020] Determining the gain may include multiplying an old gain by N. [0021] Determining the starting value may include generating a sequence of measured values from the first signal, fitting a function to the sequence of measured values, and calculating the starting value as a value of the function at an approximate start time of the polishing operation. [0022] In another aspect, a computer program product, tangibly encoded on a non-transitory computer readable media, includes instructions operable to cause a data processing apparatus to perform operations to carry out any of the above methods. [0023] In another aspect, a polishing system includes a rotatable platen to support a polishing pad, a carrier head to hold a first substrate against the polishing pad, an in-situ eddy current monitoring system including a sensor to generate a first signal depending on a thickness of a conductive layer on the substrate, and a controller configured carry out any of the above methods. [0024] In another aspect, a method of controlling polishing includes polishing a substrate at a first polishing station, during polishing of the substrate at the first polishing station, monitoring the substrate with a first eddy current monitoring system to generate a first signal, determining an ending value of the first signal for an end of polishing of the substrate at the first polishing station, determining a first temperature at the first polishing station, after polishing the substrate at the first polishing station, polishing the substrate at a second polishing station, during polishing of the substrate at the second polishing station, monitoring the substrate with a second eddy current monitoring system to generate a second signal, determining a starting value of the second signal for a start of polishing of the substrate at the second polishing station, determining a gain for the second polishing station based on the ending value, the starting value and the first temperature, for at least a portion of the second signal collected during polishing of at least one substrate at the second polishing station, calculating a third signal based on the second signal and the gain, and determining at least one of a polishing endpoint or an adjustment to a polishing parameter for the at least one substrate based on the third signal. [0025] Implementations may include one or more of the following features. [0026] Determining the gain for the second polishing station may further includes measuring a second temperature at the second polishing station. [0027] The gain may be calculated based on the resistivity of a layer being polished at first and second temperatures. [0028] [1+alpha (TE post −TE ini )] may be calculated, where TE post is the first temperature at the first polishing pad, TE ini is the second temperature at the second polishing pad, and alpha is a resistivity factor for a material of layer being polished. [0029] An ending value of the first signal for an end of polishing of the substrate at the first polishing station may be determined. [0030] Determining the ending value may include generating a first sequence of measured values from the first signal, fitting a first function to the first sequence of measured values, and calculating the ending value as a value of the function at an endpoint time for polishing at the first polishing station. [0031] A first thickness may be determined from the ending value and a calibration function relating thickness to signal strength. [0032] An adjusted thickness may be determined based on the first thickness, the first temperature and the second temperature. [0033] Determining the adjusted thickness may include multiplying the first thickness by [1+alpha (TE post −TE ini )] where TE post is the first temperature at the first polishing pad, TE ini is the second temperature at the second polishing pad, and alpha is a resistivity factor for a material of layer being polished. [0034] A desired value may be determined from the adjusted value and the calibration function. [0035] A starting value of the second signal for a start of polishing of the substrate at the second polishing station may be determined. [0036] Determining the starting value may include generating a second sequence of measured values from the second signal, fitting a second function to the second sequence of measured values, and calculating the starting value as a value of the second function at an approximate start time of polishing at the second polishing station. [0037] Determining the gain may include calculating a multiplier N according to [0000] N = ( D - K ) ( S - K ) [0038] where D is the desired value, S is the starting value, and K is a constant representing a value of the calibration function for zero thickness. [0039] The first temperature may be a temperature of a first polishing pad at the first polishing station and the second temperature may be a temperature of a second polishing pad at the second polishing station. [0040] The first temperature may be a temperature of a layer being polished at the first polishing station and the second temperature may be a temperature of the layer being polished at the second polishing station. [0041] In another aspect, a computer program product, tangibly encoded on a non-transitory computer readable media, operable to cause a data processing apparatus to perform operations to carry out any of the above methods. [0042] In another aspect, a polishing system include a first polishing station including a first platen to support a first polishing pad, a first in-situ eddy current monitoring system including a first sensor to generate a first signal depending on a thickness of a conductive layer on a substrate, and a first temperature sensor, a second polishing station including a second platen to support a second polishing pad and a second in-situ eddy current monitoring system including a second sensor to generate a second signal depending on the thickness of the conductive layer on the substrate, a carrier head to hold the substrate, and a controller configured carry out any of the above methods. [0043] Implementations may include one or more of the following advantages. Gain and offset of the monitoring system can be adjusted automatically to compensate for the parameters that can affect the eddy current signal. For example gain and offset can be adjusted for changes in the environmental conditions (e.g., temperature) or equipment parameters such as the thickness of the polishing pad. Reliability of the endpoint system to detect a desired polishing endpoint can be improved, and within-wafer and wafer-to-wafer thickness non-uniformity can be reduced. [0044] The details of one or more implementations are set forth in the accompanying drawings and the description below. Other aspects, features, and advantages will be apparent from the description and drawings, and from the claims. BRIEF DESCRIPTION OF THE DRAWINGS [0045] FIG. 1 illustrates a cross-sectional view of an example of a polishing station including an eddy current monitoring system. [0046] FIG. 2 illustrates a cross-sectional view of an example magnetic field generated by eddy current sensor. [0047] FIG. 3 illustrates a top view of an example chemical mechanical polishing station showing a path of a sensor scan across a wafer. [0048] FIG. 4 illustrates a graph of an example eddy current phase signal as a function of conductive layer thickness [0049] FIG. 5 illustrates a graph of an example trace from the eddy current monitoring system. [0050] FIG. 6 is a flow graph for the start of a polishing operation of a substrate in a polishing station. [0051] FIG. 7 is a flow graph for transferring a substrate from a first polishing station to a second polishing station. [0052] Like reference numbers and designations in the various drawings indicate like elements. DETAILED DESCRIPTION [0053] One monitoring technique for controlling a polishing operation is to use an alternating current (AC) drive signal to induce eddy currents in a conductive layer on a substrate. The induced eddy currents can be measured by an eddy current sensor in-situ during polishing to generate a signal. Assuming the outermost layer undergoing polishing is a conductive layer, then the signal from the sensor should be dependent on the thickness of the layer. [0054] Different implementations of eddy current monitoring systems may use different aspects of the signal obtained from the sensor. For example, the amplitude of the signal can be a function of the thickness of the conductive layer being polished. Additionally, a phase difference between the AC drive signal and the signal from the sensor can be a function of the thickness of the conductive layer being polished. [0055] Due to composition and assembly variations, eddy current sensors can exhibit different gains and offsets when measuring the eddy current. The eddy current can also be affected by variations in the environmental parameters, e.g., the temperature of the substrate during polishing. Run time variations such as pad wear or variations of the pressure exerted on the polishing pad (e.g., in an in-situ monitoring system) can change the distance between the eddy current sensor and the substrate and can also affect the measured eddy current signal. Therefore, calibration of the eddy current monitoring system may be performed to compensate for these variations. [0056] FIG. 1 illustrates an example of a polishing station 22 of a chemical mechanical polishing apparatus. The polishing station 22 includes a rotatable disk-shaped platen 24 on which a polishing pad 30 is situated. The platen 24 is operable to rotate about an axis 25 . For example, a motor 21 can turn a drive shaft 28 to rotate the platen 24 . The polishing pad 30 can be a two-layer polishing pad with an outer layer 34 and a softer backing layer 32 . [0057] The polishing station 22 can include a supply port or a combined supply-rinse arm 39 to dispense polishing liquid 38 , such as slurry, onto the polishing pad 30 . [0058] The carrier head 70 is operable to hold a substrate 10 against the polishing pad 30 . The carrier head 70 is suspended from a support structure 60 , e.g., a carousel or a track, and is connected by a drive shaft 74 to a carrier head rotation motor 76 so that the carrier head can rotate about an axis 71 . Optionally, the carrier head 70 can oscillate laterally, e.g., on sliders on the carousel or track 60 ; or by rotational oscillation of the carousel itself. In operation, the platen is rotated about its central axis 25 , and the carrier head is rotated about its central axis 71 and translated laterally across the top surface of the polishing pad 30 . Where there are multiple carrier heads, each carrier head 70 can have independent control of its polishing parameters, for example each carrier head can independently control the pressure applied to each respective substrate. [0059] The carrier head 70 can include a retaining ring 84 to hold the substrate. In some implementations, the retaining ring 84 may include a highly conductive portion, e.g., the carrier ring can include a thin lower plastic portion 86 that contacts the polishing pad, and a thick upper conductive portion 88 . In some implementations, the highly conductive portion is a metal, e.g., the same metal as the layer being polished, e.g., copper. [0060] A recess 26 is formed in the platen 24 , and a thin section 36 can be formed in the polishing pad 30 overlying the recess 26 . The recess 26 and thin pad section 36 can be positioned such that regardless of the translational position of the carrier head they pass beneath substrate 10 during a portion of the platen rotation. Assuming that the polishing pad 30 is a two-layer pad, the thin pad section 36 can be constructed by removing a portion of the backing layer 32 . [0061] The polishing station 22 can include a pad conditioner apparatus with a conditioning disk 31 to maintain the condition of the polishing pad. [0062] In some implementations, the polishing station 22 includes a temperature sensor 64 to monitor a temperature in the system. Although illustrated in FIG. 1 as positioned to monitor the temperature of the polishing pad 30 and/or slurry 38 on the pad 30 , the temperature sensor 64 could be positioned inside the carrier head to measure the temperature of the substrate 10 . [0063] The polishing station may include an in-situ monitoring system 40 . The in-situ monitoring system 40 generates a time-varying sequence of values that depend on the thickness of a layer on the substrate 10 . In particular, the in-situ monitoring system 40 can be an eddy current monitoring system. Similar eddy current monitoring systems are described in U.S. Pat. Nos. 6,924,641, 7,112,960 and 7,016,795, the entire disclosures of which are incorporated herein by reference. [0064] In some implementations, a polishing apparatus includes additional polishing stations. For example, a polishing apparatus can include two or three polishing stations. For example, the polishing apparatus can include a first polishing station with a first eddy current monitoring system and a second polishing station with a second eddy current monitoring system. [0065] For example, in operation, bulk polishing of the conductive layer on the substrate can be performed at the first polishing station, and polishing can be halted when a target thickness of the conductive layer remains on the substrate. The substrate is then transferred to the second polishing station, and the substrate can be polished until an underlying layer, e.g., a patterned dielectric layer. [0066] FIG. 2 illustrates a cross sectional view of an example magnetic field 48 generated by an eddy current sensor 49 . The eddy current sensor 49 can be positioned at least partially in the recess 26 (see FIG. 1 ). In some implementations, the eddy current sensor 49 includes a core 42 having two poles 42 a and 42 b and a drive coil 44 . The magnetic core 42 can receive an AC current in the drive coil 44 and can generate a magnetic field 48 between the poles 42 a and 42 b . The generated magnetic field 48 can extend through the thin pad section 36 and into the substrate 10 . A sense coil 46 generates a signal that depends on the eddy current induced in a conductive layer 12 of the substrate 10 . [0067] FIG. 3 illustrates a top view of the platen 24 . As the platen 24 rotates, the sensor 49 sweeps below the substrate 10 . By sampling the signal from the sensor at a particular frequency, the sensor 49 generates measurements at a sequence of sampling zones 96 across the substrate 10 . For each sweep, measurements at one or more of the sampling zones 96 can be selected or combined. Thus, over multiple sweeps, the selected or combined measurements provide the time-varying sequence of values. In addition, off-wafer measurements may be performed at the locations where the sensor 49 is not positioned under the substrate 10 . [0068] Referring back to FIGS. 1 and 2 , in operation, an oscillator 50 is coupled to drive coil 44 and controls drive coil 44 to generate an oscillating magnetic field 48 that extends through the body of the core 42 and into the gap between the two magnetic poles 42 a and 42 b of the core 42 . At least a portion of magnetic field 48 extends through the thin pad section 36 of the polishing pad 30 and into substrate 10 . [0069] If a conductive layer 12 , e.g., a metal layer, is present on the substrate 10 , the oscillating magnetic field 48 can generate eddy currents in the conductive layer. The generated eddy currents can be detected by the sense coil 46 . [0070] As the polishing progresses, material is removed from the conductive layer 12 , making the conductive layer 12 thinner and thus increasing the resistance of the conductive layer 12 . Therefore, the eddy current induced in the layer 12 changes as the polishing progresses. Consequently, the signal from the eddy current sensor changes as the conductive layer 12 is polished. FIG. 4 shows a graph 400 that illustrates a relationship between conductive layer thickness and the signal from the eddy current monitoring system 40 . [0071] In some implementations, the eddy current monitoring system 40 outputs a signal that is proportional to the amplitude of the current flowing in the sense coil 46 . In some implementations, the eddy current monitoring system 40 outputs a signal that is proportional to the phase difference between the current flowing in the drive coil 44 and the current flowing in the sense coil 46 . [0072] The polishing station 22 can also include a position sensor 80 , such as an optical interrupter, to sense when the eddy current sensor 49 is underneath the substrate 10 and when the eddy current sensor 49 is off the substrate. For example, the position sensor 80 can be mounted at a fixed location opposite the carrier head 70 . A flag 82 can be attached to the periphery of the platen 24 . The point of attachment and length of the flag 82 is selected so that it can signal the position sensor 80 when the core 42 sweeps underneath the substrate 10 . [0073] Alternately, the polishing station 22 can include an encoder to determine the angular position of the platen 24 . The eddy current sensor can sweep underneath the substrate with each rotation of the platen. [0074] In operation, the polishing station 22 uses the monitoring system 40 to determine when the bulk of the filler layer has been removed and/or an underlying stop layer has been exposed. The in-situ monitoring system 40 can be used to determine the amount of material removed from the surface of the substrate. [0075] Returning back to FIGS. 1 and 3 , a general purpose programmable digital computer 90 can be connected to a sensing circuitry 94 that can receive the eddy current signals. Computer 90 can be programmed to sample the eddy current signal when the substrate generally overlies the eddy current sensor 49 , to store the sampled signals, and to apply the endpoint detection logic to the stored signals and detect a polishing endpoint and/or to calculate adjustments to the polishing parameters, e.g., changes to the pressure applied by the carrier head, to improve polishing uniformity. Possible endpoint criteria for the detector logic include local minima or maxima, changes in slope, threshold values in amplitude or slope, or combinations thereof. [0076] Components of the eddy current monitoring system other than the coils and core, e.g., the oscillator 50 and sensing circuitry 94 , can be located apart from the platen 24 , and can be coupled to the components in the platen through a rotary electrical union 29 , or can be installed in the platen and communicate with the computer 90 outside the platen through the rotary electrical union 29 . [0077] In addition, computer 90 can also be programmed to measure the eddy current signal from each sweep of the eddy current sensor 49 underneath the substrate at a sampling frequency to generate a sequence of measurements for a plurality of sampling zones 96 , to calculate the radial position of each sampling zone, to divide the amplitude measurements into a plurality of radial ranges, to and to use the measurements from one or more radial ranges to determine the polishing endpoint and/or to calculate adjustments to the polishing parameter. [0078] Since the eddy current sensor 49 sweeps underneath the substrate 10 with each rotation of the platen, information on the conductive layer thickness is being accumulated in-situ and on a continuous real-time basis. During polishing, the measurements from the eddy current sensor 49 can be displayed on an output device 92 to permit an operator of the polishing station 22 to visually monitor the progress of the polishing operation. By arranging the measurements into radial ranges, the data on the conductive film thickness of each radial range can be fed into a controller (e.g., computer 90 ) to adjust the polishing pressure profile applied by a carrier head. [0079] In some implementations, the controller may use the eddy current signals to trigger a change in polishing parameters. For example, the controller may change the slurry composition. [0080] FIG. 5 shows a trace 500 generated by an eddy current monitoring system. As noted above, the signal can be sampled to generate one or more measurements 510 for each scan of the sensor across the substrate. Thus, over multiple scans, the eddy current monitoring system generates a sequence of measured values 510 . This sequence of measured values can be considered the trace 500 . In some implementations, measurements within a scan or from multiple scans can be averaged or filtered, e.g., a running average can be calculated, to generate the measurements 510 of the trace 500 . [0081] The sequence of measured values can be used to determine an endpoint or a change to the polishing parameters, e.g., to reduce within-wafer nonuniformity. For example, a function 520 (of measured value versus time) can be fit to the measured values 510 . The function 520 can be a polynomial function, e.g., a linear function. An endpoint 540 can be predicted based on a calculated time at which the linear function 520 reaches a target value 530 . [0082] As stated above, due to assembly variations and changes over time of environmental or system parameters, the eddy current monitoring system may need calibration. The eddy current monitoring system may be calibrated when the sensor is initially installed on a chemical mechanical polishing station 22 . The eddy current monitoring system may automatically be calibrated each time a substrate is loaded for polishing and/or may be calibrated during polishing. [0083] The signal from the eddy current sensor can be affected by drift of environmental parameters, e.g., temperature of the eddy current sensor itself. Drift compensation, e.g., as described in U.S. Pat. No. 7,016,795, can be performed, to compensate for some changes. However, this drift compensation technique might not address various sources of change to the signal, and might not meet increasingly stringent process demands. [0084] A measurement of the substrate from an in-line or stand-alone metrology station can be used in conjunction of measurements from the in-situ eddy current sensor to calibrate a gain of the eddy current monitoring system. For example, a desired starting signal from the in-situ eddy current sensor can determined based on the measurement from the metrology station and a calibration curve. An adjustment for the gain can then be calculated based on a comparison of the expected starting signal to the actual starting signal from the in-situ eddy current sensor. [0085] In some implementations, the calibrations can be performed using equation (1) to correct the gain. In equation (1), N is a correction factor for correcting the gain. D is the desired eddy current signal for a measured conductive layer thickness. S is the starting measured value, i.e., the measured eddy current signal at the beginning of polishing, and K is a constant representing a desired value at an off-wafer location. K can be set to a default value. [0000] N =( D−K )/( S−K )  (1) [0086] A new gain G′ can be calculated based on an old gain G and the correction factor, e.g., G′=G*N. [0087] In some implementations, the correction factor calculated from the values for S and D from one substrate is used to adjust the gain for the in-situ monitoring system for that substrate. For example, the calibration can be represented as G n =G n−1 *N n , where G n is the gain used for adjusting the n th substrate, G n−1 is the gain used for adjusting the (n−1) th substrate, and N n is the correction factor calculated from the values for S and D determined from data from the n th substrate. [0088] In some implementations, the correction factor calculated from the values for S and D from one substrate is used to adjust the gain for the in-situ monitoring system for a subsequent substrate. For example, the calibration can be represented as G n+1 =G n−1 *N n , where G n+1 is the gain used for adjusting the (n+1) th substrate, G n−1 is the gain used for adjusting the (n−1) th substrate, and N n is the correction factor calculated from the values for S and D determined from data from the n th substrate. [0089] In some implementations, the desired eddy current signal D may be calculated based on a pre-established (i.e., prior to polishing of the substrate) calibration curve relating thickness to eddy current signal. FIG. 4 illustrates an example of a calibration curve 410 . In some implementations, the calibration curve is based on eddy current signal values collected from a “golden” polishing station. As a result, ideally all polishing stations would generate the same eddy current signal for the same conductive layer thickness. [0090] FIG. 6 shows a process 600 for controlling substrate polishing, for example, chemical mechanical polishing. A measurement zone is selected on a substrate ( 610 ). The zone can be a radial range of the substrate. For example, a radial range that is empirically determined based on prior measurements to have low axial asymmetry can be selected. For example, the zone can be a radial range that excludes both the center and edge of the substrate. For example, the zone can be a radial range from 20 to 40 mm from the center of the substrate. In some implementations the zone can be selected by a user, e.g., based on input into a graphical user interface. [0091] Prior to polishing, the thickness of an outer conductive layer is measured in the selected zone ( 620 ). The outer conductive layer can be a metal layer, such as copper. The measured thickness is stored as the initial conductive layer thickness. This thickness measurement is not performed by the in-situ monitoring system. Rather, the thickness measurement may be performed by an in-line or a stand-alone metrology system suitable for measuring conductive layer thickness, such as an eddy current metrology system, e.g., the iMap™ radial scan system available from Applied Materials. [0092] The substrate is loaded to a polishing station that includes an eddy current monitoring system ( 630 ). Loading of the substrate to the polishing station may occur after the thickness of the initial conductive layer is measured. As an example, the substrate may be loaded to the polishing station 22 having an in-situ monitoring system 40 . [0093] The substrate is polished and a “raw” eddy current signal from the selected zone of the substrate is received ( 640 ). As an example, the raw eddy current signal may be received by the computer 90 of the polishing station 22 . As described above, the computer 90 may receive the raw eddy current signal for the entire substrate, and the received signal may be sampled, the position on the substrate for each sampled measurement may be determined, and the sampled measurements may be sorted into a plurality of zones including the selected zone. As noted with respect to FIG. 3 , the computer 90 may also receive the raw eddy current signal from off-wafer locations, e.g., when the eddy current sensor is not under the substrate. [0094] The received eddy current data is adjusted by a previously calculated gain and offset ( 650 ). For example, an adjusted signal value V′ can be calculated from the raw signal value V based on V′=V*G+K. [0095] In some implementations, for the n th substrate, the gain is calculated based on data from polishing a preceding (n−1) th substrate. For example, [0000] V′ n =V n *G n−1 +K and G n−1 =G n−2 *N n−1 [0000] where V′ n is the adjusted signal value for the n th substrate, V n is the raw signal value for the n th substrate, G n−1 is the gain used for adjusting the (n−1) th substrate, G n−2 is the gain used for adjusting the (n−2) th substrate, and N n−1 is the correction factor calculated from the values for S and D determined from data from the (n−1) th substrate. [0096] Gain and offset calculations for calibrating the eddy current measurement sensor are described in details with respect to step ( 670 ) below. [0097] In some implementations, when polishing a first substrate, for example, the first substrate in a batch or the first substrate after a polishing pad has been replaced, so that prior data is unreliable or unavailable, the gain is simply set as at a default value G 0 , so that [0000] V′ 1 =V 1 *G 0 +K. [0098] A new gain (or an adjustment for the gain) is calculated based on the received eddy current data and the previously measured initial conductive layer thickness ( 660 ). As an example the gain calculations can be performed by the computer 90 of the polishing station 22 . For example, a correction factor N for the gain can be calculated according to [0000] N =( D−K )/( S−K ). [0099] The initial thickness IT of the conductive layer in the selected zone was measured in step 620 . The desired value D corresponding to the initial thickness IT can be calculated from the calibration curve 410 (see FIG. 4 ). [0100] The starting measured value S can be determined from the eddy current data. That is, the eddy current data received during an initial period of polishing should correspond to the initial thickness. For example, once sufficient data has been collected during polishing, a function can be fit to the sequence of adjusted values. The function can be a polynomial function, e.g., a linear function. [0101] A value S at an initial time T 0 can be calculated from the fitted function (see FIG. 5 ). The time T 0 is not necessarily the exact start time for the polishing operation, e.g., the time that the substrate is lowered into contact with the polishing pad, but could be several seconds, e.g., 2 or 3 seconds, thereafter. Without being limited to any particular theory, using the time that the substrate is lowered into contact with the polishing pad can give an artificially high signal value since the polishing rate can initially be limited, e.g., due to the fact that the platen is still ramping up to the target rotation rate. [0102] K can be a default value. K can correspond to the value of the calibration curve 410 for zero thickness of the layer. If drift compensation is performed, e.g., as described in U.S. Pat. No. 7,016,795, then the drift compensation can automatically adjust the off-wafer signal back to K at each scan. [0103] A new gain G′ can then be calculated from an old gain G as G′=G*N. [0104] In some implementations, the new gain is used for a subsequent substrate (i.e., the substrate after the one used to generate the values for S and D). In this case, the gain to be used for the (n+1) th substrate can be expressed as G n+1 =G n−1 *N n . [0105] In some implementations, after sufficient data is accumulated to determine the starting value S for a current substrate, the new gain is calculated and a net set of data is calculated for the current substrate using the new gain. In this case, the gain to be used for the n th substrate can be expressed as G n =G n−1 *N n . [0106] For example, the new gain can be used for the current substrate when polishing a first substrate, for example, the first substrate in a batch or the first substrate after a polishing pad has been replaced. For the later polished substrates, the new gain can be used for the subsequent substrate. [0107] In some implementations, the sequence of gain values is filtered, e.g., to dampen wafer-to-wafer noise such that gain changes more smoothly, to generate a filtered gain value for a given substrate. This filtered gain value can then be used in place of G in the equations above. For example, the gain can be subject to a recursive notch filter. [0108] In either implementation, the adjusted data is used to determine a polishing endpoint or modify the polishing parameters ( 670 ). The adjusted data can represent the thickness of the conductive layer being polished and may be used to trigger a change in polishing parameters. An example of finding a polishing endpoint is described above with reference to FIG. 5 . [0109] In some implementations, one or more measurement zones may be selected and the thickness of more than one zone may be used for calibrating the eddy current sensor. In some implementations, the measurement of the initial thickness is performed at one point of the selected zone. In some implementations, the measurement is performed at two or more points of the selected zone and the measured data are averaged. [0110] The resistivity of the conductive layer may change as the temperature of the conductive layer changes. The induced eddy currents and therefore the measured eddy current signals depend on the resistivity of the conductive layer. Polishing a substrate increases the substrate temperature and reduces the induced eddy current signals. [0111] If a substrate is moved from a first in-situ polishing station to a second in-situ polishing station to continue polishing, both polishing stations using eddy current monitory, the temperature variations between the two polishing stations affects the eddy current signals. Temperature compensation can be performed as below: [0000] P post =P ini [1+alpha( TE post −TE ini )]  (2) [0000] T ini =T post [P post P ini ]  (3) [0112] In the above equations (2) and (3), P post is the resistivity factor of the layer at the second polishing station and P ini is the resistivity factor of the same layer at the first polishing station. TE post is the temperature at the second polishing station and TE ini is the temperature at the first polishing station. The parameter alpha can be calculated empirically and is a value very close to zero, e.g, 0.002 to 0.005, e.g., 0.0032. The parameter alpha can depend on the composition of the layer being polished. In some implementations, the parameter alpha can be selected by user input, e.g., the user can select from a menu listing layer compositions, and the parameter alpha corresponding to the selected layer composition is determined from a look-up table. Equation (5) can be used for thickness correction between two polishing stations having different temperatures. [0113] FIG. 7 shows a process 700 for controlling polishing when transferring a substrate from a first polishing station to a second polishing station. A measurement zone is selected on a substrate ( 710 ). As noted above, the zone can be a radial range of the substrate. For example, a radial range that is empirically determined based on prior measurements to have low axial asymmetry can be selected. In some implementations the zone can be selected by a user, e.g., based on input into a graphical user interface. [0114] The substrate is polished at a first polishing station and a “raw” eddy current raw signal from the selected zone of the substrate is received ( 720 ). As an example, the eddy current signal may be received by the computer 90 of the polishing station 22 . As described above, the computer 90 may receive the eddy current signal of the entire substrate and the sampled measurements can be sorted into different zones, including the selected zone. [0115] The received eddy current data of the selected zone of the first polishing station is adjusted by a first gain and offset ( 730 ). Gain and offset can be received from a preceding substrate measurement or may be calculated from the eddy current data of the current substrate being polished as described in details with respect to step ( 670 ) above. A first function can be fit to the eddy current data collected at the first polishing station. The first function can be a first polynomial function, e.g., a first linear function. In some implementations, for a short period of time (e.g., 10 seconds) after polishing begins the eddy current data may not be reliable and may be discarded. [0116] A first temperature of the polishing process at the first polishing station is determined ( 740 ). In some implementations, the first temperature is a temperature of the polishing pad. Alternatively or in combination, the temperature of the substrate being polished may be measured. Contacting sensors and/or non-contacting sensors (e.g., infrared sensors) may be used to measure the temperature. The temperature can be measured periodically and/or around the time polishing at the first polishing station is halted. [0117] The substrate is transferred to a second polishing station and a second temperature of the process at the second polishing station is measured ( 750 ). In general, the temperature of the same element as the first polishing station can be measured. For example, if the first temperature is a temperature of the polishing pad at the first polishing station, then the second temperature is a temperature of the polishing pad at the second polishing station. Similarly, if the first temperature is a temperature of the substrate at the first polishing station, then the second temperature is a temperature of the substrate at the second polishing station. The temperature can be measured periodically and/or around the time polishing at the second polishing station begins. [0118] Alternatively, rather than measure the second temperature at the second polishing station, the system can simply assume that the substrate is at a default temperature, e.g., room temperature, e.g., 21° C., when polishing begins at the second polishing station. [0119] The substrate is polished at the second polishing station and a raw eddy current signal the selected zone of the substrate is received ( 760 ). As an example, the eddy current signal may be received by the computer 90 of the polishing station 22 . As described above, the computer 90 may receive the eddy current raw data of the entire substrate and the sampled measurements can be sorted into different zones, including the selected zone. A second function can be fit to the eddy current data collected at the second polishing station. The second function can be a second polynomial function, e.g., a second linear function. [0120] The received eddy current data for the second polishing station is adjusted by a second gain and offset ( 770 ). Gain and offset can be received from a preceding substrate measurement or may be calculated from the eddy current data of the current wafer being polished as described in details with respect to step ( 670 ) above. In some implementations, as described in equations (2) and (3), the gain can be adjusted to incorporate the difference between the first and the second temperatures. [0121] For example, when switching the substrate from a first polishing station to a second polishing station, a correction factor N for the gain can be calculated according to [0000] N =( D′−K )/( S′−K ) [0122] The starting measured value S′ at the second polishing station can be determined from the eddy current data collected at the second polishing station. For example, once sufficient data has been collected during polishing at the second polishing station, the second function, e.g., the second linear function, is fit to the sequence of adjusted values. The value S at an initial time T 0 at the second polishing station can be calculated from the second fitted function. The time T 0 is not necessarily the exact start time for the polishing operation at the second polishing station, e.g., the time that the substrate is lowered into contact with the polishing pad, but could be several seconds, e.g., 2 or 3 seconds, thereafter. [0123] The final thickness T post of the conductive layer in the selected zone at the first polishing station can be determined. In some implementations, the first function is used to calculate a final value DF for the time TF at which polishing was actually stopped at the first polishing station. In some implementations, the final value DF is simply the target value 530 . The final thickness T post corresponding to the final value DF can be calculated based on the calibration curve 410 (see FIG. 4 ). [0124] To perform temperature compensation, an adjusted initial thickness T ini for the second polishing station is calculated based on the final thickness T post and the temperatures at the two polishing stations. For example, the adjusted initial thickness can be calculated according to T ini =T post (P post /P ini ). The desired value D's corresponding to the adjusted initial thickness T ini can then be calculated from the calibration curve 410 (see FIG. 4 ). Calculation of the gain can then proceed as discussed above. [0125] The above described polishing apparatus and methods can be applied in a variety of polishing systems. Either the polishing pad, or the carrier heads, or both can move to provide relative motion between the polishing surface and the substrate. For example, the platen may orbit rather than rotate. The polishing pad can be a circular (or some other shape) pad secured to the platen. Some aspects of the endpoint detection system may be applicable to linear polishing systems, e.g., where the polishing pad is a continuous or a reel-to-reel belt that moves linearly. The polishing layer can be a standard (for example, polyurethane with or without fillers) polishing material, a soft material, or a fixed-abrasive material. Terms of relative positioning are used; it should be understood that the polishing surface and substrate can be held in a vertical orientation or some other orientation. [0126] Embodiments can be implemented as one or more computer program products, i.e., one or more computer programs tangibly embodied in a non-transitory machine readable storage media, for execution by, or to control the operation of, data processing apparatus, e.g., a programmable processor, a computer, or multiple processors or computers. A number of embodiments of the invention have been described. Nevertheless, it will be understood that various modifications may be made without departing from the spirit and scope of the invention. For example, more or fewer calibration parameters may be used. Additionally, calibration and/or drift compensation methods may be altered. Accordingly, other embodiments are within the scope of the following claims.

A method of controlling polishing includes polishing a substrate at a first polishing station, monitoring the substrate with a first eddy current monitoring system to generate a first signal, determining an ending value of the first signal for an end of polishing of the substrate at the first polishing station, determining a first temperature at the first polishing station, polishing the substrate at a second polishing station, monitoring the substrate with a second eddy current monitoring system to generate a second signal, determining a starting value of the second signal for a start of polishing of the substrate at the second polishing station, determining a gain for the second polishing station based on the ending value, the starting value and the first temperature, and calculating a third signal based on the second signal and the gain.

This nonprovisional application claims priority to German Patent Application No. DE 10 2009 021 153.5, which was filed in Germany on May 13, 2009, and to U.S. Provisional Application No. 61/177,805, which was filed on May 13 2009, and which are both herein incorporated by reference. BACKGROUND OF THE INVENTION 1. Field of the Invention The present invention relates to a method for obtaining field strength information and to a circuit arrangement. 2. Description of the Background Art A method of this type is known from European Pat. No. EP 1 318 623 B1, which corresponds to U.S. Pat. No. 6,922,553, and which is incorporated herein by reference. Inter alia, a distance between the transmitter and receiver can be determined by means of the field strength information and can be used as a means for detecting relaying. In general, a so-called “RSSI (Received Signal Strength Indicator)” value, which is related to the input signal of a receiving antenna, is determined from the received signal strength. In this regard, an input voltage for an operational amplifier is generated from a received electromagnetic signal via a receiver unit with an input resistance, and an output voltage is provided by the operational amplifier by means of a fixed amplification factor. The input voltage at the operational amplifier is changed until the output voltage lies within a predefined interval that includes the value of the reference voltage. To this end, the input voltage at the divider node of a voltage divider is tapped and, to adjust the output voltage to the reference voltage, the resistance of the voltage divider is changed by means of connection or disconnection of resistance branches connected to one another in a complex parallel circuit. Furthermore, the method is used to determine the distance of the receiver unit from a transmitting unit from the determined field strength values. Because of the fixedly predefined values of the resistance branches of the parallel circuit, the spatial resolution is very different depending on the distance; the spatial resolution is insufficient particularly at a small distance. Furthermore, the total resistance changes depending on the selected parallel resistance. The load of an antenna circuit connected upstream changes as a result. Because of the change in load, the impedance of the antenna circuit changes in turn and in the case of an inductive coupling, feedback to the transmission circuit of a transmitter occurs in addition. Inter alia, methods for obtaining field strength information are used in systems for contactless data transmission. An important field of application is identification systems, which are used, for example, in a motor vehicle, for access control. Such identification systems include a base unit and one or more transponders; such methods are also used for determining the field strength values or the distance preferably in the transponders, which are constructed, for example, as passive and must draw the energy needed for operation by absorption from the electromagnetic field of the base station. SUMMARY OF THE INVENTION It is therefore an object of the present invention to provide a method via which a field strength of an input signal can be determined. A further object of the present invention is to provide a circuit arrangement for carrying out the method. According to a first embodiment of the present invention, a method is provided for obtaining field strength information from a received electromagnetic signal by a receiver unit, whereby an input voltage for an operational amplifier is generated from the received signal in the receiver unit by an input resistance realized as a voltage divider circuit and an output voltage is generated by the operational amplifier by a fixed amplification factor. The input voltage is changed until the output voltage lies within a predefined interval that includes the value of the reference voltage. To this end, the input voltage is tapped at the voltage divider circuit and to change the input voltage at the voltage divider circuit, which has a plurality of divider nodes and a constant resistance value, an appropriate divider node is selected and a partial voltage is tapped. As each divider node is assigned a predefined value of a first variable, from a comparison of the value, assigned to the selected divider node, of the first variable with a field strength value assigned to the value of the first variable in a value table, the field strength value received by the receiver unit can be determined. According to a second embodiment of the invention, a circuit arrangement for carrying out the method is provided, having a receiver unit for receiving an electromagnetic signal, which comprises an input resistance realized as a voltage divider circuit and an operational amplifier with a fixed amplification factor, whereby the input of the operational amplifier is connected to the voltage divider circuit and the operational amplifier supplies an output voltage from a present input voltage. Further, the receiver unit comprises a control element, connected to the output of the operational amplifier, with an input for the output voltage and an input for a reference voltage. The voltage divider circuit has a constant resistance value and comprises a plurality of divider nodes in series to one another. Further, the control element is set up, to change the input voltage, depending on a comparison of the output voltage with the reference voltage, to provide a control signal for selecting the divider node, and is set up to determine the received field strength, to provide a signal assigned to the control signal at an output, whereby the value of the signal (RS) depends on the received field strength. An advantage of the method or the circuit arrangement is that the field strength of the received electromagnetic signal can be determined, preferably without further calculations, in a simple manner from a comparison of the selected divider node and stored field strength values, assigned beforehand to the particular divider node. As a result, the method can be used for determining the spatial distance, i.e., the distance between the receiver unit and a transmitting unit, by linking the field strength values, assigned in the value table to the respective divider nodes, to a distance value. Further, by means of the values, determined in such a way, for the field strength or the distance, the data rate can be changed, for example, increased with the decline in the distance and vice versa. Further, because of the low power consumption, the method is especially suitable for obtaining field strength information in the case of passive transponders. The method of the invention is based, inter alia, on the fact that the field strength is linked to the distance between a receiver station and a transmitter station by a nonlinear functional relation. The nonlinear relation can be compensated substantially inversely proportional by the suitable realization of different resistance values of components connected between the particular divider nodes of the voltage divider circuit. In this case, the data of the connected components are accordingly determined beforehand, in order to reproduce the described functional relation between field strength and distance, so that a line as straight as possible results in a preferably linear plot of the divider node number versus the logarithm of the distance between the receiver and the base station. Because the amplification of the operational amplifier remains constant, the output voltage of the operational amplifier has a fixed ratio to the input voltage. As a result, neither at low input voltages nor at high input voltages is the amplification increased and the ratio of the signal voltage to the noise voltage remains low over the entire amplification range. Another important advantage is that over the entire detectable distance range there is an almost uniform spatial resolution both at small distances to the transmitting station and at great distances to the transmitting station. Another advantage is that the value of the input resistance of the voltage divider circuit does not change with the change in the input voltage. The voltage divider circuit functions comparably to a potentiometer with a center tap. It is understood that the input resistance of the operational amplifier is very high, so that the current, flowing from a selected node connected to the operational amplifier, is negligibly small. Overall, a change in the load at an antenna circuit, connected upstream of the receiver circuit, is avoided. In particular, during use in a transponder, the input voltage source, therefore the antenna circuit, is constantly loaded and the quality of the antenna circuit in particular is not changed. This is an important advantage in an inductively coupled antenna circuit. Tests by the applicant have shown that the voltage divider circuit, i.e., the input resistance, can be made very high-impedance. This is advantageous specifically in circuits operating passively, i.e., those that draw the needed power exclusively from the electromagnetic field of a transmitting station, in order to increase the communication range between a receiver and a transponder operating according to the method of the invention. Further, particularly with varying signal levels, a rapid turn-on and turn-off transient behavior at the output of the operational amplifier is achieved, because the output voltage and the input voltage of the operational amplifier fluctuate only to a minor extent and the capacitances, present at the input or output of the operational amplifier, are not recharged and the power consumption of the circuit unit declines. In particular, in the case of pulse-pause modulated signals, the field gaps can be reduced by the short turn-on and turn-off transient behavior and the data rate increased. Another advantage of the fixed amplification factor is that in the case of the operational amplifier, the amplification can be adjusted precisely with a low current consumption by means of a high-impedance negative feedback. An embodiment provides that the partial voltage is tapped by a controllable switching assigned to the divider node. MOS transistors are particularly suitable as the switch. Tests by the applicant have shown that the switch are controlled by an output signal of a counter component, whereby a signal, assigned to the received field strength, is provided at another output of the counter. In this case, the counter can have a number of components corresponding to the number of nodes or switches. In another embodiment, the voltage divider circuit can be a series connection of a plurality of components. The voltage divider circuit can be formed from a series connection of MOS transistors, preferably NMOS transistors, or passive resistors. In the embodiment as MOS transistors, the control inputs, i.e., the gates of the transistors, are interconnected overall. It is provided to adjust the gate voltage of the transistors in such a way that for the entire voltage range of UIN the operating point of the MOS transistors lies within the linear range of the output characteristic field and, as a result, a behavior comparable to an ohmic resistor is present. As, to connect the input of the operational amplifier to the node to be selected according to the received field strength, the voltage divider circuit and the associated switch in each case have only small capacitances and require only a very low current, great fluctuations as well in the received field strength can be detected rapidly and reliably by the method of the invention and the associated circuit arrangement and the distance can also be determined rapidly and reliably in the case of moving objects. Stated differently, the circuit arrangement, particularly because of the low capacitances, has a high switching frequency; i.e., changes in the received field strength are corrected rapidly and reliably by means of the control element and finally the received field strength is determined. According to another embodiment, it is advantageous that in each case two resistance values, compared with one another, between two divider nodes form a logarithmic relationship. As the value of the field strength in zero approximation is inversely logarithmic from the distance of the base station to the transponder, a substantially linear dependence between the number of the node and the distance can be achieved from the number of the node, directly by means of an assignment of the node number to an assigned distance value, preferably stored in a value table. In this case, it is advantageous, when the value of the signal, present at the output of the counter or of the control element, preferably indicates the number of the selected divider node. As a result, a sufficient accuracy or resolution in the distance determination can be established both at small and larger distances between the base station and the transponder. Further, the distance between base station and transponder can be determined without additional calculations by means of the values stored in the value table, also independent of the selected relationship of two resistance values. In summary, the method for determining the field strength can be easily refined into a method for determining the distance between the base station and transponder. In another embodiment, nonlinear dependences between the node numbers and the distance or a linear dependence between the node number and the field strength can be established by selection of other functional relationships between two resistance values of the voltage divider circuit. Further scope of applicability of the present invention will become apparent from the detailed description given hereinafter. However, it should be understood that the detailed description and specific examples, while indicating preferred embodiments of the invention, are given by way of illustration only, since various changes and modifications within the spirit and scope of the invention will become apparent to those skilled in the art from this detailed description. BRIEF DESCRIPTION OF THE DRAWINGS The present invention will become more fully understood from the detailed description given hereinbelow and the accompanying drawings which are given by way of illustration only, and thus, are not limitive of the present invention, and wherein: FIG. 1 shows an amplifier circuit with a voltage divider which provides a variable input voltage for an operational amplifier via a control element; FIG. 2 shows a voltage divider circuit by means of a series connection of resistors and a control circuit; FIG. 3 shows a voltage divider circuit by a series connection and a combination of a series and parallel connection of resistors; FIG. 4 shows a voltage divider circuit by a series connection of MOS transistors; and FIG. 5 shows a schematic comparison of the spatial resolution according to the method of the invention and the state of the art. DETAILED DESCRIPTION The task of the circuit arrangement, shown in FIG. 1 , is to generate an input signal UIN, which has an alternating voltage form and is proportional to the field strength of an electromagnetic signal received by a receive antenna (not shown), from an input voltage IN by means of a voltage divider circuit and to amplify an operational amplifier V 1 by a fixed factor, in order to hold the output signal UOUT within a predefined interval by means of a control element STG, which selects a suitable node and hereby does not change the resistance value of an input resistance RI, realized as a voltage divider. The input signal UIN declines completely across the input resistance RI. The voltage divider circuit with a plurality of nodes has a first terminal, which is linked to the input signal UIN, a second terminal, which is linked to the reference potential, a third terminal, which is formed as a control input and is linked to the control element STG and at which a control signal CS is present, and a fourth terminal, which is linked to an input of the operational amplifier V 1 . Further, the control element STG has an input, which is linked to the output voltage UOUT, and an input at which a reference voltage UREF is present, and an output for outputting a signal RS. The control element STG has a control unit CU and a counter ZA. The control unit CU, at which the output voltage UOUT and the reference voltage UREF are present, is linked via a control line UP and a control line DO to a counter VA, which has a plurality of components (not shown). The principle of operation will be described in greater detail below. The operational amplifier V 1 is supplied with the input voltage IN from the input signal UIN by means of the voltage divider circuit, by linking one of the nodes of the voltage divider circuit to the operational amplifier V 1 . In the case of the voltage divider circuit, a node in the vicinity of the reference potential is selected by the control element STG at high present input voltages by means of the control input and of the control signal CS, in order to lower the input voltage IN of the operational amplifier V 1 until the output voltage UOUT corresponds to a reference voltage likewise present at the input of the control element STG and lies in particular within an interval spanning the reference voltage UREF. If the output voltage UOUT is below the voltage interval, a node, which is more distant from the reference potential, is selected by the control element STG and as a result of this the input voltage IN of the operational amplifier V 1 is increased. The associated field strength value or the distance can be determined from the specifically output signal RS by reading out the quantity assigned in a memory, preferably an assigned field strength value, or most preferably an assigned distance value. Tests by the applicant have shown that the field strength value or the distance value can also be determined without the formation of a memory area with values, assigned to selected nodes, by an algorithm, preferably with use of a processor. An exemplary embodiment of the input resistance RI realized as a voltage divider circuit is shown in FIG. 2 . A series connection of individual resistors R 1 , R 2 , R 3 , to Rn- 1 , RN is formed between the first terminal of the voltage divider circuit, at which the input voltage UIN is present, and the reference potential. A node RK 0 is formed between the signal UIN and resistor R 1 , a node RK 1 between resistor R 1 and resistor R 2 , a node RK 2 between resistor R 2 and resistor R 3 , a node RK 3 after the resistor R 3 , and a node RKN- 1 between resistor RN- 1 and resistor RN. The node RK 0 can be linked by the controllable switch SM 0 to the input of the operational amplifier V 1 . To this end, the control input of the switch is linked to a component Q 0 of the counter ZA and a signal S 0 is provided by the component Q 0 . The node RK 1 can be linked by means of the controllable switch SM 1 to the input of the operational amplifier V 1 . To this end, the control input of the switch SM 1 is linked to a component Q 1 of the counter ZA and a signal S 1 is provided by the component Q 1 . The node RK 2 can be linked by means of the controllable switch SM 2 to the input of the operational amplifier V 1 . To this end, the control input of the switch SM 2 is linked to a component Q 2 of the counter ZA and a signal S 2 is provided by the component Q 2 . The node RKN- 1 can be linked by means of the controllable switch SMN- 1 to the input of the operational amplifier V 1 . To this end, the control input of the switch SMN- 1 is linked to a component QN of the counter ZA and a signal SN is provided by the component QN. If the output voltage UOUT is below the reference voltage UREF, a node above of the node selected thus far, i.e., a node that is closer to the input voltage UIN, is chosen by means of the control line DO by the control unit CU, in that, for example, the component Q 2 is deactivated by the control line DO and no corresponding signal S 2 to close the switch SM 2 is provided for the switch SM 2 and the switch SM 2 consequently breaks the connection of the node RK 2 with the input of the operational amplifier V 1 . The component Q 1 is controlled by the control unit CU by means of the control line DO and, as a result, the node RK 1 is linked to the input of the operational amplifier V 1 by means of the signal S 1 and the switch SM 1 . It is assumed that the output voltage UOUT now lies within a predefined interval and the control process is ended hereby. As the node RK 1 is selected by the control unit CU by means of the control line, a corresponding signal RS with a lower value is output by the control unit STG. A correspondingly assigned field strength value or an assigned distance value, which corresponds to the present input signal, can be read preferably out of a memory unit by means of the signal RS, and the distance between the transmitting unit and the receiver unit determined. After this, the value can be output and/or used for control of communication parameters between the transmitter and receiver unit. As the values of the individual resistors linked in series are suitably selected, the increment of the voltage changes between the individual nodes can be established in such a way that there is a sufficient spatial resolution over the entire measurable range in the particular field strength range or distance range. Another exemplary embodiment of an input resistance RI realized as a voltage divider circuit is shown in FIG. 3 . The differences to the embodiment explained in relation to the drawing documents of FIG. 2 are set forth below. A node RKM- 1 is formed between resistor RM- 1 and resistor RM. The node RKM- 1 can be linked by means of the controllable switch SMM- 1 to the input of the operational amplifier V 1 . To this end, the control input of the switch SMN- 1 is linked to a component QM (not shown) of the counter ZA and a signal SM is provided by the component QM. A resistor RP is connected between the node RKM- 1 and the node RKN- 1 . The gradation of the resistance values and thereby the voltage steps can be influenced advantageously by means of the parallel connection of the resistor RP. Particularly during tapping of the voltage at a node which is connected parallel by the resistor RP, the resistance is reduced by the parallel connection. Another embodiment of the voltage divider as a series connection of NMOS transistors is shown in FIG. 4 . Each of the transistors has an output, an input, and a control input. The control inputs of the transistors are interconnected. According to the present embodiment, the present input voltage UIN can be supplied to the input of the operational amplifier by the first switch SMTK 0 . To this end, the switch SMTK 0 is controlled by means of the signal S 0 from the component Q 0 . Over the further course, a transistor T 1 is linked to a transistor T 2 with the formation of a node TK 1 . The node TK 1 is connected by means of a controllable switch SMTK 1 to the operational amplifier V 1 (not shown) in order to supply the input voltage IN to the operational amplifier V 1 . The switch SMTK 1 has a control input, which is connected to the counter component Q 1 (not shown) and the signal S 1 is present. Transistor T 2 is linked to a transistor T 3 with the formation of a node TK 2 . The node TK 2 is connected by means of a controllable switch SMTK 2 to the operational amplifier V 1 (not shown) in order to supply the input voltage IN to the operational amplifier V 1 . The switch SMTK 2 has a control input, which is connected to the counter component Q 2 (not shown), and the signal S 2 is present. Transistor TN- 1 is linked to a transistor TN with the formation of a node TKN- 1 . The node TKN- 1 is connected by means of a controllable switch SMTKN- 1 to the operational amplifier V 1 (not shown) in order to supply the input voltage IN to the operational amplifier V 1 . The switch SMTKN- 1 has a control input, which is connected to the counter component QN (not shown) and the signal SN is present. The interconnected control inputs of transistors T 1 , T 2 , T 3 , . . . TN- 1 , TN are linked to a constant voltage source UC. The constant voltage source UC provides a control voltage UG and is linked to a supply voltage VDD and to a reference potential, preferably a ground potential. To generate the control voltage UG, the constant voltage source UC has a constant current source Iconst, which is connected in series to a transistor TR connected as a diode. The constant control voltage UG is generated by means of the constant current through the transistor diode. Tests by the applicant have shown that the control voltage UG of the transistors is selected preferably in such a way that for the entire voltage range of UIN the operating point of the MOS transistors lies within the linear range of the output characteristic field and, as a result, a behavior comparable to an ohmic resistor is present. The node whose node voltage must be supplied as the input voltage IN to the operational amplifier V 1 is selected by control of the particular switch, so that the output voltage UOUT of the operational amplifier V 1 lies at least within the required range around the reference voltage UREF. If this condition is not met, a node that meets the stated conditions is selected by the control element STG. In FIG. 5 , in a linear diagram, the distance of two control points VAL according to the present invention is plotted along the Y-axis versus the logarithm of the distance DI between a transmitting station and a base station along the X-axis; in this case, the values on both axes are selected purely arbitrarily. The relation between the distance of two divider node voltage values and the distance and hereby the spatial resolution according to the present invention is illustrated by means of a line N, whereas the course of the spatial resolution according to the state of the art, for example, known from the patent publication EP 1 318 623 B1, is shown by means of a curve P. In contrast to the prior art, according to the present invention, the determination of the distance at greater distances as well becomes possible reliably and with a much improved concentration of values. The latter is possible only very inaccurately in the prior art. In the exemplary embodiments provided thus far, the depicted NMOS transistors can also be replaced by PMOS transistors or by bipolar transistors. The control circuit of the control inputs is to be adjusted accordingly hereto. Inter alia, in an embodiment with bipolar transistors, an increased current uptake and the associated reduction of communication range, particularly in passive transponders, are a disadvantage. The invention being thus described, it will be obvious that the same may be varied in many ways. Such variations are not to be regarded as a departure from the spirit and scope of the invention, and all such modifications as would be obvious to one skilled in the art are to be included within the scope of the following claims.

A method for obtaining field strength information from a received electromagnetic signal by a receiver unit is provided, whereby an input voltage for an operational amplifier is generated from the received signal in the receiver unit by means of an input resistance realized as a voltage divider circuit and an output voltage is generated by the operational amplifier by means of a fixed amplification factor. The input voltage is changed until the output voltage lies within a predefined interval that includes the value of the reference voltage. The input voltage is tapped at the voltage divider circuit and, to change the input voltage at the voltage divider circuit, which has a plurality of divider nodes and a constant resistance value, a divider node is selected and a partial voltage is tapped. The field strength value received by the receiver unit is determined from a comparison of a quantity assigned to the selected divider node.

BACKGROUND OF THE INVENTION 1. Field of the Invention [0001] The present invention relates to a circuit and method for protecting circuit data of a programmable gate array from being copied. 2. Description of the Related Art [0002] A gate array such as an FPGA (Field Programmable Gate Array), in which a circuit configuration is rewritable by a user, is adopted to develop devices in a short period of time and mass-produce them. A circuit data used for the configuration is often stored in a storage device such as a ROM (Read Only Memory), and the FPGA is provided with SRAM (Static Random Access Memory) cells as a program device. It is necessary to copy the circuit configuration data from the ROM into the SRAM cells for programming the FPGA at every start-up such as power-on. This is because a data held in the SRAM cells is erased on power-off. For this reason, the circuit configuration data is copied from the ROM into the SRAM cells. In this case, because the ROM can easily be copied by using a ROM writer, the circuit configuration data is subject to unauthorized copying. Parts such as the FPGA and the ROM are easily commercially available. That is, a device formed by the FPGA can be developed in a short period of time but has a risk that the device developed with great resources may easily be copied. [0003] For this reason, various techniques are known for the copy protection of the device using the FPGA. A copy protection system for a programmable gate array is disclosed in Japanese Laid Open Patent Application (JP-P2003-84853 A), which is composed of a logic circuit (CPLD) and a programmable gate array circuit (FPGA). The logic circuit (CPLD) is programmed in a factory, and has an initial-state generator, a first sequence generator, and an encoding circuit. The programmable gate array circuit (FPGA) is programmed. The FPGA has a second sequence generator, a third sequence generator, a decoding circuit, and a sequence comparing circuit. The second sequence generator is a replica of the first sequence generator. In this system, an initial state is generated by the initial state generator of the CPLD. The CPLD initializes the first sequence generator to the initial state. Then, the CPLD encodes the initial state by the encoding circuit, and transmits the encoded initial state to the FPGA. The FPGA decodes the encoded initial state by the decoding circuit. The FPGA initializes the second sequence generator to the initial state. The FPGA generates a call sequence by using the third sequence generator, and transmits a call sequence to the first sequence generator and the second sequence generator. The first sequence generator generates a first response sequence based on the initial state and the call sequence. The first sequence generator transmits a first response sequence to the sequence comparing circuit. The second sequence generator generates a second response sequence based on the initial state and the call sequence, and transmits the second response sequence to the sequence comparing circuit. The sequence comparing circuit compares the first and second response sequences. The sequence comparing circuit permits an operation of an FPGA program when the first and second response sequences are identical to each other. In this way, the copy-protection system for the programmable gate array prevents the illegal copying. [0004] Also, a circuit data protecting method for a field programmable gate array provided with a volatile memory is disclosed in Japanese Laid Open Patent Application (JP-P2001-325153 A). In this method, a nonvolatile memory is provided outside the field programmable gate array, as a means to write the circuit data into the field programmable gate array. Encoded circuit data has been written in the nonvolatile memory. At the time of power-on, the circuit data is written from the nonvolatile memory into the volatile memory provided to the field programmable gate array. The field programmable gate array decodes the encoded circuit data and writes the circuit data to the volatile memory. Thus, the circuit data is protected. [0005] As described above, in these conventional examples, a copy of the circuit configuration data is prevented by using the ROM for storing the circuit configuration data for configuring the FPGA and giving a special function to the FPGA. For this reason, the conventional examples are not user-friendly. SUMMARY OF THE INVENTION [0006] In an aspect of the present invention, a copy protection circuit to prevent illegal copy of a circuit configuration data when the user circuit data is read out and transferred from a storage device to a rewritable gate array (FPGA), includes a control circuit, a data generating circuit and a data switch circuit. The control circuit controls transfer of the circuit configuration data from the storage device to the FPGA, and the data generating circuit generates pseudo circuit configuration data. The data switch circuit transfers to the FPGA, the circuit configuration data read out from the storage circuit and the pseudo circuit configuration data outputted from the data generating circuit. The data switch circuit transfers to the FPGA the circuit configuration data which is less than a data amount that the FPGA needs, and then the pseudo circuit configuration data. [0007] Here, the data generating circuit may generate the pseudo circuit configuration data based on a preset data, or may generate the pseudo circuit configuration data based on a data that has been stored in the storage circuit. Also, the pseudo circuit configuration data may contain at least a part of the circuit configuration data. [0008] Also, the storage device has stored the circuit configuration data which contains a first circuit configuration data and a second circuit configuration data. The switch circuit may transfer the first circuit configuration data read out from the storage device and then the pseudo circuit configuration data to the FPGA. The control circuit determines whether or not transfer of the second circuit configuration data to the FPGA is permitted, based on a confirmation signal outputted from a first circuit generated based on the first circuit configuration data set in the FPGA. then, the control circuit may control the switch circuit to transfer the second circuit configuration data read out from the storage circuit and then the pseudo circuit configuration data to the FPGA, based on the determination result. In this case, the copy protection circuit may further include a decoding circuit configured to decode the confirmation signal, the confirmation signal contains an encoded data; and a determining circuit configured to determine validity of the first circuit configuration data based on the decoded confirmation signal. The control circuit controls the switch to transfer the second circuit configuration data to the FPGA, when the determination result indicates validity. [0009] Also, the circuit configuration data may contain a third circuit configuration data between the first circuit configuration data and the second circuit configuration data. The circuit data generating circuit may generate the pseudo circuit configuration data based on the third circuit configuration data. [0010] Also, in another aspect of the present invention, a copy protection system includes a gate array (FPGA) that a circuit configuration is rewritable by a user; a storage device configured to store a plurality of circuit configuration data used to define the circuit configuration of the FPGA; and a copy protection circuit configured to transfer the plurality of circuit configuration data to the FPGA. The copy protection circuit includes a control circuit configured to control transfer of each of the plurality of circuit configuration data to the FPGA; and a data generating circuit configured to generate a pseudo circuit configuration data. The copy protection circuit determines whether to permit each of the plurality of circuit configuration data to be transferred to the FPGA, reads out each of the plurality of circuit configuration data from the storage device to transfer to the FPGA when the transfer is permit, and then transfers the pseudo circuit configuration data to the FPGA. [0011] Here, a first one of the plurality of circuit configuration data may have a code string. The copy protection circuit may receive the code string from the FPGA after transferring the first circuit configuration data and the pseudo circuit configuration data to the FPGA, may determine validity of the first circuit configuration data based on the received code string, and may read out a second one of the plurality of circuit configuration data to transfer to the FPGA, when the validity is confirmed. [0012] Also, the code string may be encoded and transmitted from the FPGA. [0013] Also, the storage device may include a semiconductor storage device, or a gate array programmed by a semiconductor vender. [0014] Also, in another aspect of the present invention, a method of preventing an illegal copy of a circuit configuration data for a gate array (FPGA) that a circuit configuration is rewritable by a user, is achieved by transferring a first one of a plurality of circuit configuration data stored in a storage device to the FPGA; by generating a pseudo circuit configuration data; by transferring the pseudo circuit configuration data to the FPGA after the circuit configuration data is transferred from the storage device to the FPGA; and by determining whether to transfer a second one of the plurality of circuit configuration data to the FPGA. [0015] Here, when the first circuit configuration data contains a code string, the determining step may be achieved by confirming validity of the first circuit configuration data based on the code string which is outputted from a first circuit set in the FPGA based on the first circuit configuration data; and by transferring the second circuit configuration data to the FPGA when the validity is confirmed. [0016] Also, the code string may be encoded and outputted from the first circuit. [0017] Also, the generating may be achieved by generating the pseudo circuit configuration data based on a preset data, or by generating the pseudo circuit configuration data containing at least a part of the plurality of circuit configuration data. BRIEF DESCRIPTION OF THE DRAWINGS [0018] FIG. 1 is a block diagram showing a configuration of a FPGA programming system according to an embodiment of the present invention; [0019] FIG. 2 is a diagram showing a circuit data stored in a ROM used in the embodiment of the present invention; [0020] FIGS. 3A and 3B are diagrams showing a state that the circuit data is set in the FPGA in the embodiment of the present invention; [0021] FIG. 4 is a flow chart showing an operation when the FPGA is programmed; [0022] FIG. 5 is a flow chart showing an operation when the circuit data is read out from the ROM and transferred to the FPGA; [0023] FIG. 6 is a flow chart showing another operation when the circuit data is read out from the ROM and transferred to the FPGA; and [0024] FIG. 7 is a diagram showing a start-up key contained in a start-up key code string. DESCRIPTION OF THE PREFERRED EMBODIMENTS [0025] Hereinafter, a copy protecting system of the present invention will be described in detail with reference to the attached drawings. [0026] In the present invention, a start-up control circuit data is firstly transferred from a storage device to an FPGA. After the start-up control circuit data is transferred to the FPGA, an additional data is generated in a gate array, and the additional data is transferred to the FPGA. Consequently, the data is written into an entire an SRAM storage area in the FPGA. A start-up control circuit of the FPGA is started up, and validity of a system is confirmed. Subsequently, a user circuit data is transferred from the storage device to the FPGA. After the user circuit data is transferred to the FPGA, another additional data is generated in the gate array and is then be transferred to the FPGA, and the data is written into the entire SRAM storage area in the FPGA. Thus, the validity of the system can be confirmed, and the additional data does not need to have been stored in the storage device. [0027] The copy protecting system according to an embodiment of the present invention will be described below with reference to FIGS. 1 to 4 . FIG. 1 shows a configuration of a system for programming a programmable gate array (FPGA: Field Programmable Gate Array). [0028] The FPGA programming system has an FPGA 30 in which a circuit defined by a user is programmed; a ROM 20 as a storage device which stores a circuit data to be set in the FPGA 30 ; and a gate array (GA) 10 for reading the circuit data from the ROM 20 and writing the read circuit data into the FPGA 30 . [0029] The FPGA 30 is provided with SRAM (Static Random Access Memory) as a program device. The FPGA 30 holds the circuit data in the storage area of the SRAM, and achieves a circuit operation set based on the circuit data. The circuit data is written into the SRAM from outside at the time of power-on of the FPGA programming system and so on. Usually, unless the circuit data for a storage capacitance of the SRAM is written, a start-up of a circuit cannot be carried out. That is, it is necessary to write the circuit data for the storage capacitance of the SRAM into the SRAM. [0030] In the ROM 20 , the circuit data to be set in the FPGA 30 has been stored. The circuit data is also referred to as an FPGA program. As shown in FIG. 2 , the ROM 20 stores a start-up control circuit data 21 and a user circuit data 25 . The start-up control circuit data 21 is a circuit configuration data for a start-up control circuit 31 that should be set to the FPGA 30 firstly after the power-on. The start-up control circuit 31 operates, and confirms the validity of the FPGA programming system. The start-up control circuit data 21 contains a start-up key code string 22 . On the other hand, the user circuit data 25 is a circuit configuration data for a circuit developed and designed by the user. Therefore, a data amount of the user circuit data 25 and functions thereof vary according to the user-designed circuit. Only valid circuit data of each of the start-up control circuit data 21 and the user circuit data 25 are stored in the ROM 20 . Therefore, the start-up control circuit data 21 or the user circuit data 25 alone cannot satisfy the storage capacitance of the SRAM in the FPGA 30 . Also, a boundary data 27 may be further stored between the start-up control circuit data 21 and the user circuit data 25 . The boundary data 27 is not necessarily required. Also, the boundary data 27 may not be a data for setting a valid circuit as the circuit data of the FPGA 30 . [0031] In this way, it becomes difficult to read and analyze the contents of the ROM 20 and extract only the user circuit data 25 , by providing the boundary data 27 between the start-up control circuit data 21 and the user circuit data 25 , or by continuously arranging the start-up control circuit data 21 and the user circuit data 25 . [0032] The gate array 10 is provided with a control circuit 12 , an address generating circuit 13 , a data generating circuit 14 , a data switching circuit 15 , a start-up key data 17 , a determining circuit 18 , and a decoding circuit 19 . [0033] The control circuit 12 controls each of the circuits incorporated into the gate array 10 , and controls transfer of the circuit data from the ROM 20 to the FPGA 30 . The address generating circuit 13 generates an address of the ROM 20 in response to an instruction from the control circuit 12 , and supplies the generated address ad to the ROM 20 . The data generating circuit 14 generates a pseudo circuit data in response to an instruction from the control circuit 12 , and supplies the generated pseudo circuit data to the data switching circuit 15 . The pseudo circuit data generated by the data generating circuit 14 may be a preset data, or may be generated based on the data read out from the ROM 20 . In response to an instruction of the control circuit 12 , the data switching circuit 15 selects a data to be outputted, from the data dr read out from the ROM 20 and the data outputted from the data generating circuit 14 , and outputs the selected data to the FPGA 30 . The start-up key data 17 holds a data of a start-up key for confirming the validity of the FPGA programming system, and is supplied to the determining circuit 18 . The decoding circuit 19 decodes the encoded start-up key code string outputted from the FPGA 30 , and supplies the decoded start-up key code string to the determining circuit 18 . The determining circuit 18 compares the start-up key data 17 and the decoded start-up key code string from the decoding circuit 19 , and checks whether the both are coincident with each other. The determining result is outputted to the control circuit 12 . [0034] Next, the start-up key code string used to confirm the validity of the FPGA programming system will be described. FIG. 7 shows a configuration of the start-up key code string ky 50 transmitted to confirm the validity of the FPGA programming system. The start-up key code string ky 50 contains a start-up key 54 and a data length 56 . The start-up key 54 is a preset data, and is given fixedly for every designed circuit. The start-up key 54 is stored in the gate array 10 as the start-up key data 17 , and is used to confirm the validity of the FPGA programming system. The data length 56 shows a data amount of the user circuit data 25 . [0035] Of the start-up key code string 50 , the start-up key 54 is encoded at least. The encoding method may be a private-key encoding method or a public-key encoding method. In general, the load of a decoding process is lighter in the private-key encoding method and a circuit scale is smaller, compared with the public-key encoding method. It is better to use the public-key encoding method when protecting the start-up key more strictly. [0036] Also, the start-up key 54 and the data length 56 may be encoded and transmitted and the encoded data may be transmitted with an encoded data. In this case, it is possible to reduce a risk of data leakage during transmission, during which observation is possible from outside. [0037] In addition, the data length 56 can be determined when the user circuit data 25 is stored. Therefore, the data length 56 may be stored previously in the gate array 10 without being recorded in the ROM 20 as the start-up key code string 22 . In that case, the data length 56 is not transferred between the gate array 10 and the FPGA 30 . Therefore, the observation and the data leakage at the time of the transmission do not occur. [0038] Further, the start-up key code string 22 may also be incorporated in the start-up control circuit data 21 without being encoded. In that case, an encoding circuit is incorporated into the start-up control circuit 31 , and a start-up key code string 32 is encoded to be transferred to the gate array 10 . Owing to the encoding, the observation and the data leakage at the time of the transmission do not occur. [0039] Next, an operation of the FPGA programming system for the copy protection configured as mentioned above will be described below. First, the entire operation of the gate array 10 will be described, with reference to a flow chart shown in FIG. 4 . Here, the operation will be described as the operation at the time of power-on of the FPGA programming system. [0040] However, a sequence described below may be carried out at any time when initialization of the FPGA programming system is necessary. [0041] The FPGA 30 generates a power-on detection signal if detecting the power-on, and notifies to an external unit that the reception of the circuit data is possible. Thus, the FPGA 30 enters a circuit data reception mode. On the other hand, after the power-on, the gate array 10 waits until receiving the power-on detection signal outputted from the FPGA 30 (step S 41 ). [0042] When the power-on detection signal is received and the FPGA 30 enters the circuit data reception mode (step S 41 -YES), the gate array 10 instructs the FPGA 30 to erase the circuit data set in the SRAM of the FPGA 30 (step S 42 ). This step is not necessary immediately after the power-on, if the contents of the SRAM of the FPGA 30 are all erased. However, this step is usually carried out at the time of initial setting and so on, besides the power-on. Additionally, this step is not necessary in case of the FPGA that automatically erases the circuit data stored in the SRAM when the circuit data reception mode is set. [0043] If the erasing of the circuit data in the SRAM of the FPGA 30 is completed, the gate array 10 reads the start-up control circuit data 21 from the ROM 20 , and transfers the read data to the FPGA 30 in order. Thus, the transferred data is written into the SRAM (step S 43 ). A procedure of this step is. described in detail later. [0044] The start-up control circuit data 21 is transferred to the SRAM of the FPGA 30 , and is set in the FPGA 30 , as shown in FIG. 3A . Thus, the start-up control circuit 31 is configured and the start-up key code string 22 having been stored in the ROM 20 is set as the start-up key code string 32 as it is. Then, the FPGA 30 starts up the start-up control circuit 31 (step S 45 ). [0045] The start-up control circuit 31 , if started, transfers the start-up key code string 32 incorporated therein, to the gate array 10 as the start-up key code string ky. At this time, the start-up control circuit 31 outputs the start-up key code string 32 with no change. It is preferable that the start-up key code string 32 is encoded by the start-up control circuit 31 , when the start-up key code string 32 is transferred to the gate array 10 . If the start-up key code string 32 is encoded by the start-up control circuit 31 , a value different from the start-up key code string 22 stored in the ROM 20 is generated. Consequently, a higher secrecy can be obtained. [0046] The start-up key code string 32 is transferred from the start-up control circuit 31 to the gate array 10 , as the start-up key code string ky. The start-up key code string ky 50 contains the start-up key 54 and the data length 56 , as shown in FIG. 7 . The data length 56 does not need to be transferred when the data length 56 is set previously in the gate array 10 . [0047] The gate array 10 receives the start-up key code string ky transferred from the start-up control circuit 31 (step S 47 ). The gate array 10 decodes the encoded start-up key code string ky by the decoding circuit 19 , and takes out the start-up key code string (step S 48 ). The decoded start-up key code string is compared with the start-up key data 17 by the determining circuit 18 , and the validity of the FPGA programming system is confirmed (step S 51 ). [0048] If the determining circuit 18 determines that the decoded start-up key code string is not valid (step S 51 -NG), the transfer of the user circuit data is not carried out and an FPGA start-up sequence is ended. That is, the user circuit data 25 is not transferred from the ROM 20 to the FPGA 30 . In other words, the user circuit data 25 for carrying out a desired operation cannot be set to the FPGA 30 . Additionally, the circuit data stored in the FPGA 30 may be erased. [0049] If the determining circuit 18 determines that the start-up key code string is valid (step S 51 -OK), the gate array 10 begins a transfer sequence of the user circuit data 25 . Since the circuit data of the start-up control circuit has been set in the SRAM of the FPGA 30 , the circuit data of the SRAM is firstly erased (step S 52 ). [0050] When the circuit data of the SRAM is erased, the gate array 10 transfers the user circuit data 25 of the ROM 20 to the FPGA 30 (step S 53 ). After the user circuit data 25 is all transferred to the FPGA 30 , a user circuit 35 is set to the FPGA 30 in place of the start-up control circuit 31 , as shown in FIG. 3B . The gate array 10 starts a user circuit 35 based on the user circuit data set in the FPGA 30 , and the FPGA 30 carries out the circuit operation designed by the user. [0051] In this way, it is impossible to transfer the user circuit data 25 to the FPGA 30 without the gate array 10 . Therefore, it is impossible to achieve the system function, unless the gate array 10 is owned. Therefore, the copying of the circuit data can be prevented. [0052] The transfer operation of the start-up control circuit data and the user circuit data mentioned above will be described. FIG. 5 is a flow chart showing the operation of the gate array 10 for reading the start-up control circuit data and the user circuit data from the ROM 20 , and transferring the read start-up control circuit data and user circuit data to the FPGA 30 . [0053] The control circuit 12 sets a start address of the ROM 20 in which the circuit data is stored, to the address generating circuit 13 (step S 61 ) . In case of the start-up control circuit data, it is preferable that the start address is closer to a head address of the ROM 20 than the user circuit data. The address ad of the ROM 20 is outputted from the address generating circuit 13 . The ROM 20 outputs data from the address location specified by the address ad. The gate array 10 receives the circuit data dr (step S 63 ). [0054] The data switching circuit 15 is controlled to select the data read from the ROM 20 by the control circuit 12 . The circuit data supplied to the gate array 10 is transferred to the FPGA 30 through the data switching circuit 15 . A transferred circuit data dt is stored in the SRAM of the FPGA 30 (step S 64 ). [0055] When the circuit data is transferred to the FPGA 30 , the address generating circuit 13 counts up the readout address, to instruct the ROM 20 to output the next data (step S 65 ). [0056] The steps S 63 to S 65 are repeated until the readout address exceeds a range of the circuit data (steps S 67 -NO) . Here, the end of a valid data is determined by detecting that the readout address is a last address of the circuit data. The determination of the end of the valid data may be carried out by entering a special data (pattern) to the circuit data and by detecting that special data. Also, the determination may be carried out by detecting that the transfer of the circuit data is completed so that the data is not outputted. Further, since the data amount of the circuit data is known previously based on the data length 56 , the determination may be carried out by counting the data amount by a counter and so on. [0057] When the readout address exceeds the range of the circuit data (step S 67 -YES), the control circuit 12 stops reading the circuit data from the ROM 20 , and sets the data switching circuit 15 to the side of the data generating circuit 14 . [0058] The data generating circuit 14 generates a meaningless data (such as a dummy data of all zeros and so on, which does not affect other circuits), and transfers the meaningless data to the FPGA 30 through the data switching circuit 15 . The FPGA 30 stores the meaningless data (dummy data) in the SRAM (step S 68 ). A data pattern may be set previously for the meaningless data (dummy data), or the data set in the ROM 20 may be used as the meaningless data (dummy data). Also, any circuit data stored in the ROM 20 may be used. [0059] Until the data for the storage capacitance of the SRAM of the FPGA 30 is transferred, the data generating circuit 14 generates the meaningless data (dummy data) and transfers the generated meaningless data (dummy data) to the FPGA 30 (step S 69 -NO). [0060] If the data is stored in the entire SRAM of the FPGA 30 , the configuration (writing) of the FPGA 30 is ended (step S 69 ). The termination of the circuit data transfer to the FPGA 30 is determined based on a count value of a data amount generated by the data generating circuit 14 , the count value of a transfer number by the data switching circuit 15 , or the like. [0061] In this way, the gate array 10 transfers the circuit data from the ROM 20 to the FPGA 30 . Although the address of the circuit data stored in the ROM 20 is different, the circuit data is transferred with the same procedure in both cases of the start-up control circuit and the user circuit data. Since the data generating circuit 14 generates the dummy data as padding data of the circuit data stored in the ROM 20 , the data stored in the ROM 20 may be only the circuit data valid for a circuit configuration, and the ROM capacitance can be reduced. [0062] Next, a modification of the transfer operation of the circuit data will be described with reference to FIG. 6 . FIG. 6 is a flow chart showing the operation of the gate array 10 , in which the address is outputted to the ROM 20 in the same way in the start-up control circuit data and the user circuit data, and the start-up control circuit data and the user circuit data are transferred from the ROM 20 to the FPGA 30 . [0063] The control circuit 12 sets the start address of the ROM 20 to the address generating circuit 13 (step S 71 ). In case of a single ROM, the start address of the ROM 20 is usually “ 0 ”. The same address is set as the start address of the ROM 20 in both cases of the start-up control circuit data and the user circuit data. [0064] The data read from the ROM 20 up to a position where the start-up control circuit data or the user circuit data is stored from the start address of the ROM 20 , is discarded. That is, the ROM 20 outputs the circuit data dr in response to the address outputted from the address generating circuit 13 . The gate array 10 reads the circuit data outputted from the ROM 20 (step S 73 ). The address generating circuit 13 counts up the address in response to the instruction of the control circuit 12 (step S 74 ). The control circuit 12 determines whether the address is the start address of the valid circuit data (step S 75 ). The gate array 10 repeats steps S 73 to S 75 until the start address of the valid circuit data is identified (step S 75 -NO). [0065] When the address outputted to the ROM 20 is identified as the address of the valid circuit data (step S 75 -YES), the circuit data is read out from the ROM 20 and transferred to the FPGA 30 by the gate array 10 . That is, the ROM 20 outputs the circuit data in response to the address outputted from the address generating circuit 13 . The outputted circuit data is transferred to the FPGA 30 through the data switching circuit 15 . Thus, the circuit data is stored in the SRAM (step S 77 ). At this time, the data switching circuit 15 is set to select the data read from the ROM 20 by the control circuit 12 . In order to read the next circuit data, the address generating circuit 13 counts up the address in response to the instruction of the control circuit 12 (step S 78 ). The gate array 10 repeats steps S 77 to S 79 until the address exceeds the range of the valid circuit data (step S 79 -NO). [0066] When the address exceeds the range of the valid circuit data (step S 79 -YES), the valid circuit data to be transferred to the FPGA 30 does not remain in the ROM 20 . Therefore, the gate array 10 transfers the meaningless data (which does not affect other circuits) to the FPGA 30 hereinafter such that the SRAM of the FPGA 30 is filled. On the other hand, the gate array 10 supplies the address to the ROM 20 , and operates to read the circuit data from the ROM 20 . That is, the control circuit 12 sets the data switching circuit 15 such that the data transferred to the FPGA 30 is selected from the data generating circuit 14 . The data generating circuit 14 generates the meaningless data (which does not affect other circuits), and transfers the generated meaningless data to the FPGA 30 through the data switching circuit 15 . The FPGA 30 stores the meaningless data (dummy data) in the SRAM (step S 83 ). On the other hand, if the address is supplied to the ROM 20 and the dummy data is transferred to the FPGA 30 , the address generating circuit 13 counts up the address in response to the instruction from the control circuit 12 (step S 84 ). The steps S 83 to S 85 are repeated until the data for the storage capacitance of the SRAM of the FPGA 30 is transferred (step S 85 -NO). [0067] If the data is stored in the entire SRAM area of the FPGA 30 (step S 85 -YES), a program (writing) of the FPGA 30 is ended. At this time, the address generating circuit 13 may continue to generate the address up to the final address of the ROM 20 . [0068] In this way, it is possible to transfer the circuit data without changing the address of the ROM 20 in the transfer of the start-up control circuit data and the user circuit data. Additionally, the dummy data to be transferred to the FPGA 30 is generated in the data generating circuit 14 in the above description. However, the circuit data read from the ROM 20 (the data read and discarded in the above description) may also be used. In that case, since the dummy data is the circuit data, mere observation from outside does not make it possible to distinguish the valid circuit data from an invalid circuit data. Also, when a personal computer and so on is used for a storage unit for storing the circuit data, it is preferable that the circuit data fro the start address to the end address is transferred to the gate array 10 since the circuit data is managed in a form of a file. [0069] Additionally, in the embodiment, the address generating circuit 13 is described to count up each time the data is received from the set address under the assumption that a storage device for storing the circuit data is the ROM. In case of a circuit or a device that outputs the circuit data in response to a trigger, the address generating circuit 13 may provide a signal as the trigger instead of the address. [0070] Also, it is possible to store the start-up control circuit data in an unused area other than the ROM area for storing the user circuit data and an area prepared only to write the meaningless data into the SRAM of the FPGA. In this case, it is not necessary to change the storage capacitance of the ROM even when the start-up control circuit data is needed. That is, costs of the ROM for storing the circuit data can be reduced. [0071] Also, even in a system using a generalized FPGA without providing a special function to the FPGA, and using a generalized ROM without providing a special function to the ROM, the unauthorized copying of the circuit data can be prevented. [0072] Also, according to the present invention, a storage capacitance of a storage device for storing a user circuit data can be reduced. Therefore, cost reduction is possible.

A copy protection circuit to prevent illegal copy of a circuit configuration data when the user circuit data is read out and transferred from a storage device to a rewritable gate array (FPGA), includes a control circuit, a data generating circuit and a data switch circuit. The control circuit controls transfer of the circuit configuration data from the storage device to the FPGA, and the data generating circuit generates pseudo circuit configuration data. The data switch circuit transfers to the FPGA, the circuit configuration data read out from the storage circuit and the pseudo circuit configuration data outputted from the data generating circuit. The data switch circuit transfers to the FPGA the circuit configuration data which is less than a data amount that the FPGA needs, and then the pseudo circuit configuration data.

TECHNICAL FIELD [0001] The present invention relates to a strainer filtering apparatus (referred to as a passive filtering apparatus) for filtering foreign substances, settlings, etc., generated upon occurrence of failures or accidents of an apparatus requiring a water circulation system, and more particularly, to a strainer filtering apparatus including a filtering tube used to remove foreign substances from a fluid suctioned into a pipe and a re-circulation pump when the re-circulation pump goes through an operation of an emergency core cooling system (ECCS) when a pipe failure occurs in a nuclear power plant. BACKGROUND ART [0002] A nuclear reactor of a nuclear power plant is surrounded by a safety vessel formed of concrete and steel, which is referred to as a containment, in which a coolant circulates to maintain a proper temperature. In addition, the nuclear reactor includes an ECCS for cooling the nuclear reactor upon occurrence of failures or accidents. [0003] The ECCS must be operated upon occurrence of accidents such as coolant leakage, etc., to cool the nuclear reactor for 30 days with no external interference. The ECCS is a system for collecting coolant discharged and water sprinkled upon a pipe failure into a sump disposed at the lowermost part in the containment, sprinkling the water from an upper part of the containment using the re-circulation pump to cool the containment, and circulating some of the water through a nuclear reactor cooling system to remove remaining heat of the nuclear reactor using a remaining heat removing pump. [0004] When coolant leakage occurs due to damage to a pipe, etc., in a primary system of the nuclear power plant, foreign substances such as lagging materials, coating materials, latent foreign substances, etc., are generated due to discharge of a coolant. [0005] In addition, the discharged coolant and water sprinkled from a sprinkler system of the containment move all foreign substances to a re-circulation sump disposed at a lower end of the containment of the nuclear reactor. Therefore, in order for the foreign substances not to decrease performance of the ECCS, a filtering apparatus is provided in front of an inlet part of a suction pipe guided to an emergency cooling pump. [0006] When a high temperature and high pressure pipe is broken, foreign substances such as fragments of lagging materials, coating materials, etc., are generated and moved toward the sump, and the filtering apparatus functions to filter the foreign substances moved to the sump and supply the filtered water into the re-circulation pump, without interfering with the operation of the re-circulation pump. [0007] The filtering apparatus ensures that the foreign substances generated due to accidents can be filtered and the water can appropriately pass therethrough. In this case, a pressure drop due to the foreign substances must be guaranteed not to exceed an allowable critical value. [0008] A conventional filter screen used in a pressurized water reactor type nuclear power plant has a small screen surface only, and the screen surface is mainly formed of flat grid segments. Thus, when the screen surface is contaminated with fiber settlings, a pressure drop at the screen may be largely increased to an unallowable level. [0009] However, the filtering apparatus having a single surface may be easily deformed by a high pressure, and a small effective filtering area per a unit volume may decrease filtering efficiency. In order to solve the problem, while the number of filtering apparatus may be increased, their installation cost is high, which causes economical problems. Therefore, a filtering apparatus capable of increasing a filtering area per unit volume is still needed. SUMMARY OF THE INVENTION [0010] In order to solve the foregoing and/or other problems, it is an aspect of the present invention to provide a strainer filtering apparatus including a filtering tube capable of providing a substantially larger effective filtering area in the same length and width, substantially reducing foreign substances covering a suction surface and a flow resistance of the foreign substances, and reducing a pressure drop at a cooling water pass corresponding thereto. [0011] It is another aspect of the present invention to provide a strainer filtering apparatus including a filtering tube capable of reducing manufacturing and installation costs to solve economical problems in exchange and installation thereof, rapidly manufacturing the apparatus by assembling a relatively small number of components, and maximizing a filtering area per unit volume even in a narrow space. [0012] The foregoing and/or other aspects of the present invention may be achieved by providing a strainer filtering apparatus including at least one inlet side into which cooling water is introduced and an outlet side through which the cooling water is discharged, including: a plurality of filtering tubes formed in a hollow shape by bending a punched plate having a plurality of filtering holes; an upper plate having first grooves formed at a lower surface to be coupled to upper ends of the filtering tubes and an inlet part into which the cooling water is introduced; and a lower plate having punched holes into which lower ends of the filtering tubes are coupled, wherein the filtered cooling water in the filtering tubes is introduced through the punched holes to be discharged to the outlet side. [0013] The first grooves may be formed in plural, and the punched holes may be formed in plural at positions corresponding to the first grooves so that the plurality of filtering tubes are coupled between the upper plate and the lower plate. [0014] The cooling water may be introduced into a space between the upper plate and the lower plate to contact outer surfaces of the plurality of filtering tubes. [0015] The cooling water contacting the outer surfaces may be filtered to be introduced into discharge cams in the filtering tubes. [0016] The plurality of first grooves may form a first groove arrangement group in which grooves are spaced a predetermined distance from a center of the upper plate and spaced a predetermined interval from each other, and the plurality of punched holes may be formed to correspond to the first grooves. [0017] The first groove arrangement group may be formed on the upper plate in plural, and the upper ends of the filtering tubes may be press-fitted into the first grooves. [0018] The punched holes may have a diameter equal to an inner diameter of the filtering tubes, and the lower plate may further include second grooves formed around the punched holes and equal to an outer diameter of the filtering tubes, whereby the lower ends of the filtering tubes are press-fitted into the second grooves of the lower plate. [0019] The strainer filtering apparatus may further include a coupling member installed between the upper and lower plates to couple the upper and lower plates to fix the filtering tubes between the upper and lower plates. [0020] The coupling member may include at least one fixing pin installed in a space between the upper and lower plates and fastening members for fastening both ends of the fixing pin to the upper and lower plates. [0021] The strainer filtering apparatus may further include a fixing member installed at one side of the lower plate and coupling the lower plate to the passage through which the cooling water flows. [0022] The fixing member may be provided around the lower plate in plural. The filtering holes may have a diameter of 1 to 3 mm. [0023] According to a strainer filtering apparatus of the present invention, it is possible to provide a substantially larger effective filtering area in the same length and width. Therefore, a flow resistance of settlings and foreign substances covering a suction surface can be substantially reduced. In addition, a pressure drop generated along the strainer filtering apparatus can be reduced depending on reduction in flow resistance. [0024] Further, since the strainer filtering apparatus of the present invention is fabricated by assembling a filtering tube formed of a punched plate, an upper plate and a lower plate, without welding, it is possible to easily perform maintenance and installation thereof. Furthermore, since a plurality of filtering tubes formed of a punched outer surface are vertically disposed, a load pressure can be distributed to increase structural integrity. [0025] In addition, it is possible to provide the strainer filtering apparatus capable of being rapidly assembled with a relatively small number of components, and maximizing a filtering area per unit volume even in a narrow space. BRIEF DESCRIPTION OF DRAWINGS [0026] The above and other aspects and advantages of the present invention will become apparent and more readily appreciated from the following description of exemplary embodiments, taken in conjunction with the accompanying drawings of which: [0027] FIG. 1 is a perspective view of a strainer filtering apparatus in accordance with an exemplary embodiment of the present invention; [0028] FIG. 2 is a bottom view of the strainer filtering apparatus in accordance with an exemplary embodiment of the present invention; [0029] FIG. 3 is a plan view of the strainer filtering apparatus in accordance with an exemplary embodiment of the present invention; [0030] FIG. 4 is a perspective view of a filtering tube in accordance with an exemplary embodiment of the present invention; [0031] FIG. 5 is a perspective view of an upper plate in accordance with an exemplary embodiment of the present invention; [0032] FIG. 6 is a bottom view of the upper plate in accordance with an exemplary embodiment of the present invention; [0033] FIG. 7 is a plan view of the upper plate in accordance with an exemplary embodiment of the present invention; [0034] FIG. 8 is a perspective view of a lower plate in accordance with an exemplary embodiment of the present invention; [0035] FIG. 9 is a bottom view of the lower plate in accordance with an exemplary embodiment of the present invention; [0036] FIG. 10 is a plan view of the lower plate in accordance with an exemplary embodiment of the present invention; and [0037] FIG. 11 is an exploded perspective view of the strainer filtering apparatus in accordance with an exemplary embodiment of the present invention. DETAILED DESCRIPTION [0038] Various embodiments will now be described more fully with reference to the accompanying drawings in which some embodiments are shown. These inventive concepts may, however, be embodied in different forms and should not be construed as limited to the embodiments set forth herein. Rather, these embodiments are provided so that this disclosure is thorough and complete and fully conveys the inventive concept to those skilled in the art. [0039] In the drawings, like reference numerals designate like elements throughout the invention. [0040] Hereinafter, constitution and structure of a strainer filtering apparatus 10 in accordance with an exemplary embodiment of the present invention will be described. [0041] FIG. 1 is a perspective view of a strainer filtering apparatus 10 in accordance with an exemplary embodiment of the present invention, FIG. 2 is a bottom view of the strainer filtering apparatus 10 in accordance with an exemplary embodiment of the present invention, and FIG. 3 is a plan view of the strainer filtering apparatus 10 in accordance with an exemplary embodiment of the present invention. [0042] As shown in FIG. 1 , the strainer filtering apparatus 10 in accordance with the present invention includes a plurality of filtering tubes 100 formed of a punched surface. In addition, the strainer filtering apparatus 10 includes an upper plate 200 having first grooves 210 and inlet parts 220 formed in a lower surface, a lower plate 300 having a plurality of punched holes 310 , second grooves 320 formed in an upper surface and fixing members 330 formed at a periphery, and coupling members 400 for coupling the upper plate 200 and the lower plate 300 . [0043] FIG. 4 is a perspective view of a filtering tube 100 in accordance with an exemplary embodiment of the present invention. [0044] As shown in FIG. 4 , the filtering tube 100 is hollow and formed of a punched surface. A punched plate having a plurality of filtering holes 110 is bent to form the filtering tube 100 . The filtering holes 110 formed in the filtering tube 100 may have a diameter of about 1 to 3 mm, preferably, 2 to 2.5 mm. In this embodiment, the filtering tube 100 has an outer diameter of 40 mm and an inner diameter of 36 mm. A discharge cam 120 is formed in an inner space of the filtering tube 100 to discharge the filtered cooling water to the outlet side. The strainer filtering apparatus 10 in accordance with an exemplary embodiment of the present invention includes a plurality of filtering tubes 100 . [0045] FIG. 5 is a perspective view of an upper plate 200 in accordance with an exemplary embodiment of the present invention, FIG. 6 is a bottom view of the upper plate 200 in accordance with an exemplary embodiment of the present invention, and FIG. 7 is a plan view of the upper plate 200 in accordance with an exemplary embodiment of the present invention. [0046] As shown in FIGS. 5 to 7 , the upper plate 200 includes inlet parts 220 through which cooling water is introduced. The inlet parts 220 may have various shapes, regardless of the shapes shown in FIGS. 5 to 7 . In addition, the upper plate 200 includes fixing holes 230 to be coupled to the lower plate 300 by coupling members 400 . [0047] As shown in FIG. 6 , the upper plate 200 has a plurality of first grooves 210 formed in a lower surface thereof. Upper ends of the filtering tubes 100 are press-fitted into the plurality of first grooves 210 . Therefore, the first grooves 210 have a diameter equal to an outer diameter of the filtering tubes 100 . In addition, the number of the first grooves 210 formed in the lower surface of the upper plate 200 is equal to the number of the filtering tubes 100 installed in the strainer filtering apparatus 10 . In this embodiment, the number of the first grooves 210 is 24. [0048] As shown in FIG. 6 , the first grooves 210 have an arrangement group in which the first grooves 210 are spaced a predetermined distance from a center of the upper plate 210 and spaced apart from each other at predetermined intervals. In this embodiment, the first grooves 210 have two arrangement groups. In addition, the first grooves 210 are formed in the lower surface, in which the fixing holes 230 and the inlet parts 220 are not disposed. The diameter and number of the first grooves 210 , the shape of the arrangement groups and number of the first grooves 210 will be understood not to be limited to the specific embodiment. [0049] FIG. 8 is a perspective view of a lower plate 300 in accordance with an exemplary embodiment of the present invention, FIG. 9 is a bottom view of the lower plate 300 in accordance with an exemplary embodiment of the present invention, and FIG. 10 is a plan view of the lower plate 300 in accordance with an exemplary embodiment of the present invention. [0050] As shown in FIGS. 8 to 10 , the lower plate 300 includes a plurality of punched holes 310 . In addition, second grooves 320 are formed around the punched holes 320 of an upper surface of the lower plate 300 . The punched holes 310 of the lower plate 300 have a diameter equal to an inner diameter of the filtering tubes 100 , and the second grooves 320 have a diameter equal to an outer diameter of the filtering tubes 100 . [0051] Therefore, lower ends of the plurality of filtering tubes 100 may be press-fitted into the second grooves 320 of the lower plate 300 . The punched holes 310 and the second grooves 320 of the lower plate 300 are arranged to correspond to the first grooves 210 of the upper plate 200 . That is, the first grooves 210 of the upper plate 200 and the punched holes 310 and the second grooves 320 of the lower plate 300 are formed at symmetrical positions. [0052] In addition, the lower plate 300 also includes fixing holes 230 like the upper plate 200 . As shown in FIGS. 8 to 10 , it will be appreciated that the fixing holes 230 formed in the lower plate 300 also correspond to the fixing holes 230 formed in the upper plate 200 . [0053] Further, the lower plate 300 includes fixing members 330 . As shown in FIGS. 8 to 10 , the lower plate 300 includes four fixing members 330 formed around the lower plate 300 and spaced a predetermined interval from each other. The lower plate 300 is fixedly installed at one side of the passage, through which cooling water flows, by the fixing members 330 . [0054] FIG. 11 is an exploded perspective view of the strainer filtering apparatus in accordance with an exemplary embodiment of the present invention. [0055] As shown in FIG. 11 , upper ends of the plurality of filtering tubes 100 are press-fitted into the first grooves 210 of the upper plate 200 , respectively. On the other hand, lower ends of the plurality of filtering tubes 100 are press-fitted into the second grooves 320 of the lower plate 300 , respectively. [0056] In addition, the coupling members 400 are installed between the upper plate 200 and the lower plate 300 of the strainer filtering apparatus 10 . The coupling members 400 couple the upper plate 200 and the lower plate 300 while maintaining a gap between the upper plate 200 and the lower plate 300 . Specifically, the coupling members 400 include fixing pins 410 and fastening members 420 . The fixing pins 410 have threads formed at both ends thereof. In addition, both ends of the fixing pins 410 are inserted into the fixing holes 230 of the upper plate 200 and the fixing holes 230 of the lower plate 300 , and fixed to the upper plate 200 and the lower plate 300 by the fastening members 420 such as nuts, etc. Therefore, the coupling members 400 couple the upper plate 200 and the lower plate 300 and fix the filtering tubes 100 between the upper plate 200 and the lower plate 300 . [0057] According to the above constitution, the cooling water introduced into the inlet part 220 of the upper plate 200 is introduced into a space between the upper plate 200 and the lower plate 300 to contact outer surfaces of the plurality of filtering tubes 100 , and the cooling water contacting the outer surfaces of the filtering tubes 100 is introduced into the discharge cams 120 in the filtering tubes 100 to be filtered. Then, the filtered cooling water is discharged from the discharge cams 120 of the filtering tubes 100 to the outlet side through the punched holes 310 of the lower plate 300 . [0058] The strainer filtering apparatus 10 including a filtering tube in accordance with an exemplary embodiment of the present invention can effectively increase a filtering area even under internal conditions of a narrow containment. In addition, since the apparatus is designed as a modular structure constituted by the filtering tubes 100 , the upper plate 200 , the lower plate 300 , and the coupling members 400 , installation and maintenance thereof is very easy. Further, in comparison with a conventional filtering apparatus having a single surface, it is possible to minimize deformation even under a high pressure and increase a filtering area per unit volume, securing safety of a recirculation operation of an ECCS when a pipe failure occurs in a nuclear power plant. Furthermore, the present invention can be applied to all pressurized water reactor type and pressurized heavy water reactor type nuclear power plants. [0059] The foregoing description concerns an exemplary embodiment of the invention, is intended to be illustrative, and should not be construed as limiting the invention. The present teachings can be readily applied to other types of devices and apparatus. Many alternatives, modifications, and variations within the scope and spirit of the present invention will be apparent to those skilled in the art.

A strainer filtering apparatus including a filtering tube providing a substantially larger effective filtering area for its length and width, substantially reducing foreign substances covering a suction surface and flow resistance of the foreign substances, and reducing pressure drop at a cooling water passage. The strainer filtering apparatus includes at least one inlet side into which cooling water is introduced and an outlet side through which the cooling water is discharged, hollow filtering tubes with filtering holes along their lengths, an upper plate having first grooves located at a lower surface and coupled to upper ends of the filtering tubes and an inlet part into which the cooling water is introduced, and a lower plate having punched holes to which lower ends of the filtering tubes are coupled. The cooling water in the filtering tubes is introduced through the punched holes and discharged at the outlet side.

BACKGROUND OF THE INVENTION This invention relates to a fluid valve coupling for hose or pipe ends. Many fluid valve couplings have been devised, most of which require male and female coupling members. There have, however, been devised fluid valve couplings comprising coupling members of identical construction, such as disclosed in my previously granted U.S. Pat. No. 3,217,746 and in U.S. Pat. Nos. 3,168,335; 3,176,717 and 3,564,918. SUMMARY OF THE INVENTION The present invention is a quick disconnect fluid valve coupling of basically simple design, providing a free flow of fluid under high pressure, the coupling embodying coupling members of identical construction which may be quickly and easily interengaged and retained in locked position, and which provides an efficient and tight mechanical seal, each coupling member automatically sealing itself upon disconnection. The coupling comprises a plurality of readily manufactured and assembled parts which are capable of operating within a broad range of temperatures and pressures, with caustic or non-caustic materials, the coupling members being connected by an interlocking lug system of a type which permits engagement and disengagement thereof by hand. The components of the interlocking lug system are internal of the outermost portion of the coupling to prevent the entry of foreign matter into the coupling at the locus of connection of the coupling members. The coupling members include interengageable retaining elements for maintaining the lugs of the coupling members in interlocked relationship, thereby preventing accidental disengagement of the coupling due to vibration, etc. The coupling members include breakaway locking means adapted to normally retain the coupling members in locked, coupled position, the coupling members separating upon exertion of predetermined axial stress on the coupling. The retaining elements for retaining the lugs in interlocked relationship include circular spring clips held in retaining recesses on the outer surfaces of the lugs, which clips engage complemental recesses in the outer surfaces of the lugs of the opposed coupling member. Locking means for normally preventing accidental disengagement of the coupling members include a collar around the housing of each member for urging the circular clip and complemental recesses into locking relationship. DESCRIPTION OF FIGURES OF THE DRAWINGS FIG. 1 is a plan view of the breakaway fluid valve coupling of the present invention, showing the coupling members in coupled position; FIG. 2 is an enlarged longitudinal sectional view of the breakaway fluid valve coupling of the present invention, showing the coupling members in uncoupled position; FIG. 3 is a view similar to FIG. 2, showing the coupling members in coupled position; FIG. 4 is a sectional view taken along the line 4--4 of FIG. 2, looking in the direction of the arrows; FIG. 5 is an enlarged end elevational view of one of the coupling members of the present breakaway fluid valve coupling, showing to advantage the valve sealing member and interlocking lug assembly forming a part thereof; FIG. 6 is an enlarged fragmentary sectional view of the coupling members showing to advantage the valve sealing members in the coupled position illustrated in FIG. 3; FIGS. 6a to 6e are diagrammatic views illustrating the position assumed by the sealing members during the steps of coupling of the coupling members; and FIGS. 7 to 11 are enlarged fragmentary sectional views illustrating the steps of coupling of the coupling members comprising the breakaway fluid valve coupling of the present invention. DESCRIPTION OF THE INVENTION Referring now in greater detail to the drawings, and more particularly to FIGS. 1 to 3, the breakaway fluid valve coupling of the present invention comprises a pair of coupling members designated 20 and 20', which are of identical construction. Because of the identify in construction, only coupling member 20 will be described in detail, the corresponding parts of coupling member 20' being identified in the drawings by like, primed numbers. Coupling member 20 includes a valve housing 22, a valve assembly 24 and an interlocking lug assembly 26. Valve housing 22 comprises a cylindrical body 28 which is of uniform external diameter throughout a portion of its length, one terminal of which tapers inwardly at 30 and issues into an internally threaded extension 32 of reduced diameter adapted for engagement with a hose or pipe A through which fluid passes. A portion of the internal periphery of body portion 28 is thickened to provide an abutment or a shoulder 36. An enlarged axial bore 38 is thereby formed within the body portion 28, the diameter of which bore is gradually decreased through tapering portion 30, as indicated at 40, and issues into an axial bore 42 of very slightly increasing diameter, formed within threaded extension 32. Valve assembly 24 includes a cylindrical valve body 44 having an axial bore, one end of which bore is flared to provide a tapered valve seat 26. The valve assembly bore is in communication with enlarged bore 38 of valve housing 26, and is enlarged at 48, thereby forming an internal annular shoulder or abutment 50. The terminal of cylindrical valve body 44 adjacent enlarged bore 48 is externally thickened to provide a peripheral annular abutment 52. The opposite terminal of cylindrical valve body 44 is threadedly engaged with the inner periphery of cylindrical body 34 as indicated at 54. Adjacent the point of threaded engagement of cylindrical body 44 with body portion 34, there is provided an annular recess adapted for the reception of a sealing ring 56. Cylindrical valve body 44 is also provided with an inner peripheral annular recess which is also adapted for the reception of a sealing ring 58. Valve assembly 24 further includes a movable poppet valve 60 extending longitudinally through the valve assembly bore and enlarged bore 48. The outer terminal of poppet valve 62 is exteriorly chamfered for purposes which will be hereinafter more fully set out. The opposite terminal of poppet valve 60 is gradually enlarged to form a tapered head 62, the peripheral portion of which is complemental to, and adapted for engagement with, tapered valve seat 46. Tapered head 62 is provided with an annular recess in which is a sealing ring 64 which engages tapered valve seat 46 when the valve is in the closed position shown in FIG. 2. Poppet valve 60 has an axial bore 66 which extends from the inner terminal of the valve to a point lying within tapered head 62, at which point there are provided ports 68 which, as shown to advantage in FIG. 3, permit the passage of fluid from bore 66 into the valve assembly bore, or vice versa. Ports 68 extend through a substantial portion of the poppet valve's periphery, thereby forming webs 70 connecting the main portion of valve body 60 to tapered head 62. The free terminal of tapered head 62 is cone-shaped as indicated at 72, for optimum flow characteristics of the fluid through the coupling member. As shown to advantage in FIG. 2, that portion of poppet valve 60 remote from tapered head 62 extends beyond the terminal of valve body 44 through a bore 74 formed by interlocking lug assembly 26. An end sealing member 76 of disc shape, has a central bore fitted over the end portion of poppet valve body 60 and is uniplanar with the terminal thereof. End sealing member 76 is of substantially the same diameter as the inner periphery of interlocking lug assembly 26. The central bore of sealing member 76 is enlarged at its outer extremity so that, in conjunction with the terminal outer periphery of poppet valve body 60, an annular recess 77 is formed in which a sealing ring 78, preferably an "O" ring made of buna, neoprene or other suitable material, is positioned. A convolute spring 80 lies circumjacent valve body 60, one end of which spring engages the rear face of end sealing member 76 and the other end of which engages abutment 50 of valve body 44. Spring 80 serves to effect longitudinal movement of valve assembly 24 to the extended position shown in FIG. 2, poppet valve 60 being in sealing engagement with valve seat 46. When coupling members 20 and 20' are interengaged, end sealing member 76 is forced rearwardly against the tension of spring 80 until the rear face of the end sealing member 76 engages the terminal of valve body 44 and tapered head 62 is lifted out of engagement with valve seat 46 to permit flow of fluid between ports 68 and axial bore 38. In accordance with the present invention, lug assembly 26 includes a cylindrical body portion 82, having an axial bore which is reduced in diameter from the inner terminal of the body to a point intermediate its length to provide an annulus 84 extending from the forward terminal of cylindrical body 28 to peripheral abutment 52. The inner periphery of body portion 82 is in contiguous engagement with the outer periphery of cylindrical body 44, the body portion being rotatable with respect to, and independent of, the valve body. The outer periphery of body portion 82 is extended at the terminal adjacent cylindrical body 28 to provide an annular abutment 86, the function of which will be hereinafter more fully set out. As shown in FIGS. 7 to 11, the terminal of body portion 82 remote from abutment 86 is extended on diametrically opposed sides thereof to provide transversely arcuate, rectangular, locking fingers or lugs 88 extending beyond end sealing member 76 a predetermined distance, the terminal portions of which lugs are chamfered at 90. The outer surfaces of lugs 88 are provided with pairs of spaced transverse notches or recesses including an inner recess 92 and an outer recess 94, the recesses of one lug being vertically aligned with the corresponding recesses of the other lug. The side walls defining inner recess 92 are preferably vertical, while those of outer recess 94 are inwardly sloping from top to bottom as indicated at 96, for purposes hereinafter more fully set out. A circular clip 98 of spring wire or other suitable material, lies circumjacent lugs 88, diametrically opposed portions thereof being fitted in inner recesses 92 of lugs 88. Lug assembly 26 further includes a cylindrical clip-retaining collar 100, the aft terminal of which is contiguous with the outer periphery of annular abutment 86. An internal shoulder 102 is provided adjacent the forward extremity of collar 100, the leading edge of which shoulder is sloped or beveled at 104 for exerting constant pressure on clip 98 in a direction to retain the latter within recesses 92. Pressure is maintained on the clip by a convolute spring 106 disposed in a peripheral recess 108 lying between collar 100 and cylindrical body portion 82. One terminal of spring 106 engages fixed abutment 86 and the opposite terminal thereof engages terminal shoulder 102, thereby effecting longitudinal movement of clip-retaining collar 100 to constantly maintain pressure on clip 98. Arcuate lug-receiving recesses 110, shown in FIG. 4, are formed between lugs 88 in the space separating cylindrical valve body 44 and clip-retaining collar 100. Recesses 110 are of a size and shape to receive the opposed lugs 88' of coupling member 20'. OPERATION In use of the present breakaway coupling, pipe terminals A and A' are threaded or secured in any other suitable fashion to extensions 32, 32' of valve housings 22, 22'. Lug assemblies 26, 26' are then rotated until lugs 88, 88' lie in planes at right angles to each other, as shown in FIG. 2. Coupling of the members is then effected by direct longitudinal movement of lugs 88 into lug receiving recesses 110' and lugs 88' into lug receiving recesses 110. This effects engagement of sealing rings 78, 78' of end sealing members 76, 76', thereby effecting the rearward movement of valve assemblies 24, 24' to move the latter to the operative position shown in FIG. 3. Movement of lugs 88, 88' into their respective recesses is continued until clip 98 operatively engages outer arcuate recesses 94' of lug 88', and clip 98' operatively engages outer arcuate recesses 94 of lug 88, thereby locking the coupling members in operative coupled position. The operations involved in coupling and uncoupling the breakaway coupling, can be better understood by reference to the enlarged fragmentary views of the interlocking cycle illustrated in FIGS. 7 to 11. FIG. 7 shows lug 88 initially in the process of being inserted into its mating recess 110', at which time spring clip 98' is in tension and has been lifted from retaining recess 92 by the upward force from the lateral motion imparted by the chamfered portion 90 of lug 88. Spring clip 98 is lifted out of its retaining recess 92 in similar fashion. This action creates rearward pressure on collars 100, 100', against the tension of springs 106, 106'. Continued movement of lug 88 into the recess 110' exerts a maximum force on spring clip 98 in the position of FIG. 8, which also sets up forces downwardly and forwardly because of the mechanical actions of the respective coupling parts in relation to each other. Further continued actual movement as shown in FIG. 9 places spring clip 98' at the edge of outer recess 94, this action occurring simultaneously at all lug locations. Referring to FIG. 10, it will be seen that as spring clip 98' attempts to return to its partially relaxed position against the inclined side wall of recess 94, a force is created in relation to lug 88 which is a function of a sine of angular recess 94. Also, at the same time, still another force comes into play, which is that force created by convolute spring 106' which forces beveled shoulder 104' of pressure collar 100' into engagement with spring clip 98'. By reason of spring 106' and the mechanical advantage of beveled shoulder 104', spring clip 98' is forced further into recess 94 and lug 88 is urged into recess 110'. It will be noted from a consideration of FIGS. 7 to 11 that certain relationships are important for optimum operation of the breakaway coupling. The distance between lugs 88 and the forward terminal of collar 100', as viewed in FIG. 8, is sufficient to maintain clearance for spring clip 98'; recesses 92, 92' are deep enough to prevent escape of spring clips 98, 98' therefrom (see FIG. 11); the outermost wall defining arcuate recesses is slightly lower than the inner wall, to facilitate engagement of clips 98, 98' therein; the preloaded position of spring clips 98, 98' is such that these units are forced into the position shown in FIG. 10 and the clips will lie in spaced relation to the bottom of annular recesses 94, 94' as shown in FIG. 11; the urging of spring 106, 106' is of sufficient force to maintain constant pressure on collars 100, 100'; and that the opposed faces of the coupling members never come into direct contact with each other, i.e., a clearance is maintained between these two surfaces. The mechanical advantage imposed upon lug 88 in a lefthand direction into the relation of recess 110' in a righthand direction, is in the order of 2:1 with respect to each other, in the position illustrated in FIG. 10. It is a salient feature of the present invention to provide a coupling so designed that the immediate sealing of members 76, 76' is effected by the longitudinal force exerted by convolute springs 80, 80' on sealing rings 78, 78' while simultaneously effecting the interlocking cycle illustrated in FIGS. 7 to 11. By virtue of the present arrangement of parts, the exertion of the longitudinal forces on the opposed "O" rings and the interlocking cycle occur very nearly simultaneously. The passage of fluid through the coupling is, however, effected just prior to the completion of the interlocking of the coupling members. As soon as sufficient longitudinal force is exerted against sealing rings 78, 78', the pressure exerted by convolute springs 80, 80' is overcome, thus forcing the coupling members' poppet valves open, allowing the passage of fluid immediately. Any pressures present in line A, A' at the moment of connection of the coupling members enhances the effectiveness of the sealing components, as shown in FIGS. 6a to 6e. Referring to FIGS. 6a to 6e, it will be seen that, upon interlocking of coupling members 20, 20', opposed sealing rings 78, 78' are flattened within the angular recesses 77, 77' to positively preclude any leakage of fluid therebetween. The fluid may then pass, for example, through the pipe or hose A, through axial bores 42, 40 and 38 to ports 68 of the valve assembly and thence into axial bore 66 of poppet valve 60. Fluid flow is continued through axial bore 66' of poppet valve 60', thence through ports 68' into bores 38', 40' and 42' and into connecting hose or pipe A'. It will be noted from a consideration of FIGS. 2 and 3, that valve housing 22 and valve assembly 24 are so constructed to provide a minimum of resistance to the flow of fluid through the coupling. In FIGS. 6a to 6e, there is illustrated the action of the sealing members during the coupling operation. FIG. 6a shows an end seal "O" ring within angular recess 77' prior to engagement with an opposed "O" ring. For optimum results, the "O" ring is in engagement with the side walls and bottom wall of the recess and approximately 10 percent of its cross sectional area protrudes from the recess. FIG. 6b shows the "O" ring in contact as the sealing action is initiated, as indicated by the arrows showing the direction of forces being exerted. Since the "O" rings initially protrude from the recesses by approximately 10 percent, there is at this state of the coupling action, a gap of 0.020 inches between opposed faces of sealing members 76, 76'. FIG. 6c represents the condition of the sealing "O" ring immediately prior to exertion of spring pressure generated by forcing springs 80, 80' rearwardly. Complete sealing is effected at this time and annular bores 66, 66' are in communication. However, bores 38, 40, 42 and 38', 40', 42' are not in communication until further lateral pressure required to force the convolute springs 80, 80' rearwardly against the tension thereof is exerted, and force tapered heads 62, 62' out of communication with their valve seats 46, 46'. Fluid flow through the coupling is thereby permitted. FIG. 6d represents the position of the sealing "O" rings when the coupling members are fully interlocked as shown in FIG. 3. In this position, the gap between sealing members 76, 76' has been reduced to its operational tolerance and "O" rings 78, 78' are completely captivated in their respective recesses. Thus, the rings, under pressure, cannot escape or be forced into poppet valve bores 66, 66' since they are retained by the outer wall of the poppet valves 60, 60'. Under compression, the chamfered leading terminals 61, 61' of poppet valves 60, 60' permit portions of the ring to be displaced into the cavity formed by the communication of the opposed chambered terminals. Any captured air in the areas of the sealing ring recesses is pressurized and tends to force the rings into a tighter sealing engagement with each other and their surrounding walls. FIG. 6e diametrically depicts the positions of the "O" ring when the coupling is fully connected and flow of fluid is taking place under pressure. During exertion of high pressures through bores 66, 66', these pressures are also being exerted on sealing rings 78, 78' since the pressurized fluid is attempting to pass the "O" ring seals. However, under sufficient pressure through bores 66, 66', the displaced "O" rings in the cavity formed by the leading chamfered terminals 61, 61' of poppets 60, 60', shown in FIG. 6d, now tend to be displaced further into their own restricting recesses. The pressures exerted then become a factor of the sealing forces as indicated by the arrows of FIG. 6e. The high pressures generated force "O" rings 78, 78' into the slight gap between sealing members 76, 76'. However, since the interlocking mechanism is, by design, of sufficient strength to withstand the forces exerted within the operating range, the "O" rings cannot be extruded into the gap between the sealing members without stressing the lugs out of alignment. To disconnect the coupled units by manual force, it is first necessary to move pressure collars 100, 100' in a rearward direction against tension exerted by springs 106, 106'. Spring clips 98, 98' are thereby freed to move upwardly along the sloping side walls of recesses 94, 94' in respect to the mechanical tension of spring clips 98, 98'. Spring clips 98, 98' are retained within recesses 92, 92' in the manner set out supra. In use of the breakaway coupling, it is normally required that the units be disconnected automatically, rather than manually, under predetermined force being exerted on the coupling. In automatic uncoupling, the reverse procedure would be followed than that set out above in connection with coupling the members, and a predetermined disconnection force may be designed into the coupling itself as, for example, by varying the tension of springs 106, 106'. The present breakaway coupling is also of value for automatically disconnecting and sealing fluid lines where, for example, the pressure in the fluid system suddenly increases beyond safe limits, or if external mechanical forces are exerted on the fluid lines beyond safe limits. The interlocking lug assemblies 26, 26' retain the coupling members in interlocked relationship under forces sufficient to prevent accidental disengagement due to vibrations or the like. However, upon exertion of predetermined forces acting in opposite directions, clips 98, 98' will become disengaged from outer arcuate recesses 94, 94' respectively, to effect longitudinal movement of valve assemblies 24, 24' to the inoperative position under the urging of springs 80, 80'. The breakaway coupling is particularly adapted for use in application where the pipe or hose members A, A' are normally not left in proximity to each other after uncoupling, such as of trailer tractors, where the trailers are unhooked from the cab. The use of the breakaway coupling positively precludes damage to air or hydraulic lines in the event that the coupling members are inadvertently left in coupled position during the removal of the cab from the trailer. As above set out, the forces exerted in opposite directions on the coupling members overcome the forces retaining the coupling members together, thereby effecting disconnection of the coupling members and sealing of the valve assemblies. The fluid valve coupling of the present invention is of basically simple design and comprising components requiring no special or exotic machine processing or materials. By virtue of the identity of the coupling members, a minimum number of simple parts are required for the valve coupling. The valve operation of the coupling is a straight line motion on the longitudinal axis, accomplished by normal coupling action in accordance with the procedure above set out, with no rotational forces exerted on the component parts by virtue of a freely rotating interlocking lug system. The efficient and tight mechanical seal effected by the present arrangement of parts permits use thereof in both the coupled and uncoupled state under conditions of high pressure and vacuum and operating within a broad temperature range. While there has been herein shown and described the presently preferred embodiment of this invention, it is to be understood that such has been done for purposes of illustration only and that various changes be made therein within the scope of the claims hereto appended.

A breakaway fluid valve coupling including a pair of identical coupling members each of which is connected to an end of a fluid conveying pipe. The coupling members are joined together to allow fluid to pass through the pipe and fluid valve coupling. The coupling members are pulled apart by exertion of axial stress on the members, the coupling members automatically sealing the ends to prevent leakage of the fluid.

FIELD OF INVENTION The present invention relates to a plant promoter and an application thereof. DESCRIPTION OF RELATED ART For expression of an exogeneous structural gene of interest in a plant or plant cells, cauliflower mosaic virus 35S promoter (hereinafter referred to as “35S promoter”) which consists of about 0.8 kb has been used, but the minimal region of the 35S promoter (e.g., −90 region of the 35S promoter which consists of 98 nucleotide bases (−90 to +8)) has not been satisfactory for practical expression of a gene of interest because of the low transcription activity of the region (Odell et al., Nature 313: 810-812 (1985), Jensen et al., Nature 321: 669-674 (1986), Jefferson et al., EMBO J. 6: 3901-3907 (1987), Kay et al., Science 236: 1299-1302 (1987), Sanderset et al., Nucl. Acid. Res. 4: 1543-1558 (1987), etc.), Benfey et al., EMBO J. 8: 2195-2202 (1989). SUMMARY OF INVENTION To facilitate efficient expression of an exogeneous structural gene of interest in a plant or plant cells, the present inventors have extensively studied and found that a compact specific nucleotide sequence can enhance the expression of a gene of interest when connected upstream to a promoter including the minimal region as described above. An object of the present invention is to provide a promoter functional in plant cells, which contains the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2. In addition to the specific nucleotides shown in SEQ ID NO:1 or SEQ ID NO:2, the promoter may include additional nucleotide on the 5′ and/or 3′ end of the sequence. Another object of the present invention is to provide a plasmid comprising a promoter functional in plant cells which contains the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2, and a terminator functional in plant cells. An object of the present invention is to provide a plasmid comprising a promoter functional in plant cells which contains the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2, a structural gene of interest and a terminator functional in plant cells. A further object of the invention is to provide a plasmid pGbox10 shown in FIG. 4 or pGbox11 shown in FIG. 5 . Further objects of the present invention are to provide a plant or plant cell expressing a structural gene of interest under the control of the promoter described above, and a plant harboring the plasmid described above. Still further objects of the present invention are to provide a method for expressing in plant cells a structural gene of interest, wherein the expression is controlled by the promoter described above, and a method for constructing a plasmid, which comprises connecting the promoter of the present invention, a structural gene of interest and a terminator functional in plant cells in this order. DESCRIPTION OF FIGURES FIG. 1 shows the construction of plasmid −90/GUS from pBI101. FIG. 2 shows the construction of the plasmid Gbox10/−90/GUS from plasmid −90/GUS. FIG. 3 shows the construction of the plasmid Gbox11/−90/GUS from plasmid −90/GUS. FIG. 4 shows the construction of pGbox10 from the plasmid Gbox10/−90/GUS. FIG. 5 shows the construction of pGbox11 from the plasmid Gbox11/−90/GUS. DESCRIPTION OF PREFERRED EMBODIMENT According to the present invention efficient expression of a gene of interest can be attained by the compact promoter. The genetic engineering techniques to be used in the present invention follow standard procedures such as described in J. Sambrook, E. F., Frisch, T. Maniatis “Molecular Cloning, 2nd ed.”, publ. Cold Spring Harbor Laboratory (1989), D. M. Glover, “DNA Cloning”, publ. IRL (1985), and elsewhere. First, description will be made on the promoter functional in plant cells, which contains the nucleotide sequence shown in Sequence ID No:1 or No:2. The promoter usually contains the sequence shown in SEQ ID NO:1 or SEQ. ID. NO:2, and preferably it contains multiple copies, particularly in the form of tandem repeats. The promoter preferably contains four or more such copies. The nucleotide sequence shown in SEQ ID NO: 1 or SEQ ID NO:2 may be of either natural or synthetic origin. Synthesis of the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2 may be accomplished by standard DNA chemical synthetic techniques. The present promoter usually contains a minimal element necessary for transcription initiation other than the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2. The element for transcription refers to a sequence region necessary for transcription initiation such as a transcription initiation site only, a transcription initiation site and TATA sequence, a transcription initiation site and CAAT sequence, alternatively a transcription initiation site and TATA sequence and CAAT sequence, etc. A typical example of such a sequence region is provided by SEQ ID NO:3 and includes a sequence in which the 5′-terminus corresponds to a nucleotide sequence located at least about 30 nucleotides upstream of the transcription initiation site and in which the 3′-terminus sequence corresponds to the region from the transcription initiation site to the translation initiation site. The transcription activity of such a basic region is generally low. Specified examples of such regions include the 98 nucleotide base region of the 35S promoter which includes the transcription initiation site (+8 to −90) (hereinafter referred to as the “−90 region”), the −204 to +8 region of the tomato gene encoding the small subunit of the Ribulose-1,5-diphosphate carboxylase-oxidase (rbcS-3A) (Plant Cell 1: 217-227 (1989)), the −287 to +29 region of the PR1a gene promoter (Plant Cell 2: 357-366 (1990)), and the −195 to +32 region of the potato protease inhibitor gene (PI-II) (Plant Cell 2: 61-70 (1990)). The regions described above that contain a minimal element necessary for transcription initiation are usually utilized by inserting them downstream of the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2. The nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2 and the sequence containing a minimal element for transcription initiation may be prepared by restriction enzyme digestion of genomic DNA whose sequence is known, or by Polymerase Chain Reaction (PCR) amplification of a region of a DNA nucleotide containing SEQ ID NO:1 or SEQ ID NO:2, or the sequence containing a minimal element for transcription initiation using genomic DNA as a template and appropriate oligonucleotides as primers, or by DNA chemical synthesis. When the −90 region of the 35S promoter is utilized, the promoter of the present invention may be constructed by synthesizing and annealing oligonucleotides representing the +strand (refer to Sequence No:4) and the −strand (refer to SEQ ID NO:5 of the −90 region. With respect to the 35S promoter, a deletion promoter which is shorter than the −90 region may be obtained, and is usually constructed by digestion with restriction enzymes, by PCR amplification using genomic DNA as a template and oligonucleotide primers having sequences appropriate to amplify the deletion promoter region, or by DNA chemical synthesis methods. The terminator functional in plant cells to be used in the present invention includes, for example, plant-derived terminators such as the terminator from the T-DNA-derived nopaline synthase gene (NOS), or virus-derived terminators usually used in plant genetic engineering techniques, such as the terminators from the Garlic virus GV1 and GV2 genes. The plasmid of the present invention comprises a promoter functional in plant cells, which contains the nucleotide sequence shown in SEQ ID NO:1 or NO:2, and a terminator functional in plant cells. The plasmid of the present invention is preferably constructed so as to contain at least one cloning site upstream of the terminator and downstream of the promoter to accommodate the desired structural gene. More preferably the plasmid is constructed in such a way as to contain multiple cloning sites. The term “cloning site” here refers to a region of DNA which can be recognized and digested by restriction enzymes usually utilized in genetic engineering procedures. One example of such a plasmid containing cloning sites is the plasmid pGbox10 shown in FIG. 4 or the plasmid pGbox11 shown in FIG. 5 . Examples of useful exogenous structural genes which may be expressed under the control of the present plasmid are plant defense genes such as the phenylalanine ammonia lyase gene (PAL), the chalcone synthetase gene (CHS), the chitinase gene (CHT), the lysozyme gene, the PR protein gene, etc., disease resistance genes such as Pto, and genes which increase the resistance against bacteria, fungi, viruses and insects in all plant tissues such as the virus coat protein gene, BT ( Bacillus thuringiensis ) toxin protein gene, etc. Other useful genes include those which increase the protein content of feed crops such as genes encoding storage proteins, the soybean glycinin gene, and the β-conglycinin gene, etc., genes which increase the methionine and lysine content of feed crops, such as the Brazil nut 2S albumin gene, the 10 kDa or 15 kDa protein genes from maize and rice, etc. and genes which increase the biotin content of feed crops such as the bacterial genes from Escherichia coli, etc. which encode the bioA, bioB, bioC, bioD, bioF, and bioH enzymes involved in the synthesis of biotin. Other useful genes include those which improve the quality of lipids by providing stability to oxidation, decreasing the phospholipid content and increasing the oleic acid and linoleic acid content such as the stearoyl-ACP-desaturase, acyl-ACP-thioesterase, and 3-phosphate acyltransferase genes, genes which increase the resistance to low temperatures by increasing the proportion of unsaturated fatty acids such as the acetyltransferase gene, and genes which make possible the generation of herbicide-resistant crops such as by the expression of the gene encoding L-phosphinothricin acetyltransferase, 5-enolpyrvil-3-phosphosikimate synthase or of other genes related to herbicide resistance. After the plants of the present invention are grown, the whole plants or parts of the plants can be harvested and sold. Components of the plants, especially components containing expression products of the structural gene of interest, can be separated from, extracted from and/or concentrated from the plants by conventional techniques. Likewise, plant cells can be harvested and sold as a commercial product per se (such as a food source of food additives) and/or components of the plant cells, in particular components containing expression products of the structural genes of interest, can be separated from, extracted from and/or concentrated from the plant cells from the plant cells by conventional techniques. The plasmid of the present invention may be constructed, for example, by the following method. The promoter functional in plant cells, which contains the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2, is inserted into the multicloning site of the plasmid containing the terminator functional in plant cells such as pBI101 (Clontech Inc.; Jefferson et al., EMBO J. 6: 3901-3907 (1987)). Furthermore an exogenous marker gene such as β-glucuronidase can be excised and replaced with a desired structural gene as necessary. Another possible method is to use a binary vector such as pBIN19 (Nuc. Acid. Res. 12: 8711-8721 (1984)) and insert the promoter, the desired structural gene if necessary, and the terminator to be used in the present invention, in that order, into the multicloning site. As for the methods for introducing the plasmid into plant cells, there are conventional methods such as the Agrobacterium infection method (i.e. infection of plant tissue with the soil bacteria Agrobacterium), electric-based introduction methods (electric-based method of introduction into protoplasts: electroporation), or direct introduction by a particle gun (direct introduction into plant tissues or cultured cells: particle gun method). Plant cells harboring the plasmid of the present invention may be obtained by the conventional plant tissue culture techniques described in, e.g. S. B. Gelvin, R. A. Schilperoot and D. P. S. Verma, Plant Molecular Biology/Manual, Kluwer Academic Publishers press, 1988 and subsequently plants or their parts derived from these plant cells may be obtained by regeneration according to the protocols described therein. Moreover, the method described herein may be applied for the plant species which include monocots such as rice, maize, barley, wheat, and onion, dicots such as the members of the Leguminosae, i.e. soybeans, peas, kidney beans, alfalfa, members of the Solanaceae such as tobacco, tomato, and potato, members of the Cruciferae such as cabbage, rapeseed, mustard plant, members of the Cucurbitaceae (gourd family) such as melon, pumpkin, cucumbers, members of the Ammiaceae such as carrot and celery, and members of the Compositae such as lettuce. According to the present invention, plant cells (transformed plant cells) and plants (transformed plants) expressing a structural gene of interest can be obtained efficiently by using the compact promoter or plasmid of the present invention. EXAMPLES The present invention will be further illustrated in more detail by the following examples. Example 1 Construction of GUS Expression Plasmids Gbox10/−90/GUS. Gbox11/−90/GUS and −90/GUS The plasmid −90/GUS was constructed by inserting the −90 region prepared by DNA chemical synthesis into the multicloning site upstream of the GUS gene contained within the commercially available plasmid pBI101 (Clontech Inc.; Jefferson et al., EMBO J. 6: 3901-3907 (1987)) which also contains a terminator functional in plant cells. The plasmid Gbox10/−90/GUS and plasmid Gbox11/−90/GUS were further constructed by inserting upstream of the −90 region of −90/GUS a multimer containing 4 tandem repeats of the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2, respectively. The method of construction is explained in detail below (refer to FIGS. 1, 2 and 3 ). Step 1: Synthesis and Purification of Complementary Oligonucleotides Containing 4 Tandem Repeats of G-box10 or G-box11 and Complementary Oligonucleotides Containing the −90 Region Two complementary oligonucleotides (+strand and −strand, 46 bases each, refer to SEQ ID NO:6, 7, 8 and 9) containing a multimer consisting of tandem repeats of the oligonucleotide sequence shown in SEQ ID NO:1 or NO:2 and having a HindIII restriction site at the 5′terminus and XbaI site at the 3′terminus when annealed were chemically synthesized. Further two complementary oligonucleotides (+strand and −strand, 102 bases each, refer to SEQ ID NO:4 and SEQ ID NO:5) containing the −90 region which includes the minimal elements necessary for transcription initiation and BamHI restriction sites at both termini when annealed were chemically synthesized. These chemically synthesized oligonucleotides were deprotected with ammonia treatment (55° C., 5 hr) and subsequently purified by reverse-phase HPLC (YMC GEL ODS S-5). The solvent used in purification was 0.1 M triethylamine (TEAA) and the oligonucleotides were extracted using a concentration gradient of 5-100% methyl cyanide (CH 3 CN). The extracted oligonucleotides were recovered, the residue dried, redissolved in 80% acetic acid (3 ml) for 20 min and subsequently dried again. The dried residue was redissolved in distilled water (3 ml) and redried, and this procedure was repeated for a total of three times before the sample was finally dissolved in 100 μl of TE (10 mM Tris-HCl pH 8.0, 1 mM EDTA). The oligonucleotide sample thus obtained was further purified by electrophoresis with a 5% polyacrylamide gel (80 V, 1 hr). The pair of complementary oligonucleotides containing the multimer consisting of 4 tandem repeats of the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2 (42 bases each) as well as the pair of complementary oligonucleotides containing the −90 region (102 bases each) were each extracted from the polyacrylamide gel and recovered by electrophoretic transfer (180 V, 40 mins.) out of the gel in dialysis tubing (SPECTRA/POR, molecularporous membrane tubing MW 3500). Step 2: Annealing of Complementary Strands 0.5 μg of a pair of complementary oligonucleotides containing the multimer consisting of 4 tandem repeats of the nucleotide sequence shown in SEQ ID NO:1 or NO:2, and a pair of complementary oligonucleotides containing the −90 region (102 bases) were heated to 100° C. for 3 min in 10 μl aqueous solution, transferred to a 65° C. water bath, allowed to slowly return to room temperature and finally chilled quickly in ice water. By this procedure a DNA fragment containing a multimer consisting of 4 tandem repeats of the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2, and a DNA fragment containing the −90 region (102 bases) were prepared. Step 3: Construction of the −90/GUS Plasmid 2 μg of pBI101 were digested with 10 units of BamHI and 0.1 μg of the −90 region (102 bases) prepared in step 2 and 0.5 μg of the BamHI-digested pBI101 were mixed and ligated using T4 DNA ligase (DNA ligation kit, Takara Shuzo Ltd.). This mixture was used to transform the E. coli strain HB101 (Takara Shuzo Ltd.) according to the protocol of Cohen et al. (Proc. Natl. Acad. Sci. USA, 69: 2110-2114 (1972)). From the resistant colonies grown up on LB agar plates containing 50 μg/ml kanamycin plasmid DNAs were extracted by the alkaline-SDS method and their structures were analyzed by restriction enzyme digestion. Isolated plasmids which yielded approx. 100 bp DNA fragment upon digestion with SmaI and XbaI were selected. The sequences of the DNA inserted in these selected plasmids were determined by the dideoxy method (Sanger et al., Proc. NAtl. Acad. Sci., USA, 74: 5463 (1977)) and clones in which the −90 region had been inserted in the proper orientation were chosen where the proper orientation means the same orientation of the −90 region as is found in the 35S promoter. Thus the −90/GUS plasmid was constructed (see FIG. 1 ). Step 4: Construction of the Gbox10/−90/GUS Plasmid and Gbox11/−90/GUS Plasmid 2 μg of the −90/GUS plasmid prepared in step 3 were digested with 10 units of XbaI and HindIII (37° C., 2 hrs.), and the digestion products were fractionated by electrophoresis in 0.8% low melting-point agarose (80 V, 1.5 hr). The HindIII-XbaI fragment was recovered using a centrifugation tube with a DNA recovery filter (Takara Shuzo Ltd.) and purified. 0.1 μg of the DNA fragment prepared in Step 2 containing multimer consisting of 4 tandem repeats of the nucleotide sequence shown in SEQ ID NO:1 or SEQ ID NO:2 respectively, and 0.5 μg of the HindIII-XbaI-digested −90/GUS plasmid were mixed and ligated using T4 DNA ligase (DNA ligation kit, Takara Shuzo Ltd.). This mixture was used to transform the E. coli strain HB101 (Takara Shuzo Ltd.) according to the protocol of Cohen et al. (Proc. Natl. Acad. Sci. USA, 69: 2110-2114 (1972)). From the resistant colonies grown up on LB agar plates containing 50 μg/ml kanamycin plasmid DNAs were extracted by the alkaline-SDS method and their structures were analyzed by restriction enzyme digestion. Isolated plasmids which yielded approx. 50 bp DNA fragments upon digestion with HindIII and XbaI were selected. The structures of the plasmids were confirmed by the dideoxy method (Sanger et al., Proc. Natl. Acad. Sci., USA, 74: 5463 (1977)), and the sequences of the DNA inserts were determined. Thus, the Gbox10/−90/GUS plasmid and Gbox11/−90/GUS plasmid were constructed respectively (see FIGS. 2 and 3 ). Example 2 Construction of pGbox10 and pGbox11 2 μg of the plasmid Gbox10/−90/GUS and plasmid Gbox11/−90/GUS were each digested with 10 units of SmaI (37° C., 2 hrs.) respectively. 0.1 μg of the digestion product was mixed with 0.05 μg of commercially available SacI linkers (5′-CGAGCTCG-3′) (Takara Shuzo Ltd.) and ligated using T4 DNA ligase (DNA ligation kit, Takara Shuzo Ltd.). The ligation products were subsequently digested with 10 units of Sac1, religated with T4 DNA ligase (DNA ligation kit (Takara Shuzo)) and used to transform the E. coli strain HB101 (Takara Shuzo) according to the method of Cohen et al. (Proc. Natl. Acad. Sci. USA. 69: 2110-2114 (1972)). From the resistant colonies grown up on LB agar plates containing 50 μg/ml kanamycin plasmid DNAs were extracted by the alkaline-SDS method and their structures were analyzed by restriction enzyme digestion. Isolated plasmids were subjected to SmaI and SacI digestion and single linear DNA fragments were selected. Thus the plasmid pGbox10 and plasmid pGbox11 were constructed respectively(see also FIGS. 4 and 5 ). Example 3 Preparation of Plasmid DNA for Gene Introduction The plasmids Gbox10/−90/GUS, Gbox11/−90/Gus, and −90/GUS prepared in Example 1 were further purified by cesium chloride density gradient centrifugation technique respectively. These DNA plasmids were purified by adding 1 g of cesium chloride and 80 μl of 10 mg/ml ethidium bromide per 1 ml of DNA solution. Sealing the resulting solution in centrifuge tubes (Quick-Seal, Beckman Co.) and subjecting the solution to centrifugation in an NVT65 rotor at 60,000 rpm for 24 hr, yielded the purified plasmid. Example 4 Generation of Transformed Tobacco Plant by Indirect Introduction Method Each purified plasmid described in Example 3 was introduced by heat treatment (37° C., 5 mins.) into Agrobacterium ( Agrobacterium tumefaciens LBA4404; streptomycin resistant, rifampicin resistant; Hoekma et al., Nature 303: 179-180 (1983)) treated with 20 mM CaCl 2 to make them competent. Transformants were obtained by utilizing the kanamycin resistance conferred by the plasmid NPTII gene (Trien-Cuot et al., Gene 23: 331-341 (1983)) and being selected on L agar plates containing 300 μg/ml streptomycin, 100 μg/ml rifampicin and 100 μg/ml kanamycin. The Agrobacterium transformants thus obtained were cultured at 20° C. for a day in L broth media containing 300 μg/ml streptomycin, 100 μg/ml rifampicin and 100 μg/ml kanamycin and this bacterial suspension was then used to infect a disc of tobacco plant according to the standard protocol described in S. B. Gelvin, R. A. Schilperoot and D. P. S. Verma, Plant Molecular Biology Manual, publ. Kluwer Academic Publishers, 1988. The infected tobacco leaf discs (SR-1) were cultured for 4 days in MS-NB agar media and then transferred to MS-NB agar media containing 500 μg/ml cefotaxime to kill the Agrobacteria. 11 days later the leaf discs were transferred to MS-NB agar media containing 500 μg/ml cefotaxime and 100 μg/ml kanamycin thus the selection of transformed plants was initiated. Approximately 4 weeks later young plants from which green stems and leaves had developed were separated from the leaf disc and further cultured on MS agar media containing 500 μg/ml cefotaxime and 50 μg/ml kanamycin, after which young plants giving rise to roots were selected. The selected young plants were transferred to soil and cultivated in a greenhouse to produce transformed plants. Example 5 Generation of Transformed Carrot by Indirect Introduction G-box10/−90/GUS purified in Example 3 was introduced into Agrobacterium tumefaciens LBA4404 according to the method described in Example 4 and the obtained Agrobacterium was infected to hypocotyl of a carrot species, Nantes Scarlet. The infected carrot hypocotyl was placed on LS-D agar medium containing 500 μg/ml cefotaxime. Selection of the transformed plant was initiated after 10 days later by placing the plant to LS-D agar media containing 100 μg/ml cefotaxime and 50 μg/ml kanamycin. The generated callus after one month was transferred to a LS-D agar media containing 100 μg/ml cefotaxime and 50 μg/ml kanamycin repeatedly every four weeks. The callus selected after 2 months was transferred to LS agar media and a whole plant was regenerated by way of adventitive embryo. Example 6 Production of Transformed Plant by Direct Introduction G-box10/−90/GUS purified as described in Example 3 was introduced into immature embryo of rice, Notohikari with a particle gun (Reebock company) according to a method described in Shimada, T., et al., Bull.RIAR, Ishikawa Agr. Coll.4:1-8 (1995) or Ko Shimamoto and Kiyotaka Okada, Experimental Protocol of model plant, Rice and Arabidopsis, Syujyun-sya, 1996 (ISBN4-87962-157-9 C3345). After sterilized seeds were cultured for 7 to 10 days in LS-D2 media, the embryo was picked up and placed on LS-D2 agar media with its scutellum organ upside. 10 μg of either G-box10/−90/GUS, −90/GUS, or pBI121 (35S/GUS) as a control, and pDM302 containing equimolar amount of bialaphos resistance gene (J. Cao et al., Plant Cell Rep. 11:586-591, 1992) were coated on 3 mg of a gold particle (avarage diameter 1 μm). Two shots were injected into rice embryos with each shot containing 0.2 mg of gold particle (1 μg DNA, Injection pressure 150-200 kg/cm 2 , at 70 mmHg). Two days after the introduction resulting embryos were transferred to LS-D2 agar media containing 4 mg/ml of bialaphos to select herbicide resistant cells. The thus obtained herbicide resistant callus were transferred to culturing media for redifferentiation and regenerated adventitious embryos or small plants were transferred to LS agar media containing 4 mg/ml of bialaphos to grow the herbicide resistant plants. Example 7 Confirmation of Insertion of the Introduced Gene in Transformed Plants 1. Preparation of Genomic DNA from Transformed Plants Genomic DNA was isolated from transformed plants according to the CTAB method described in Hirofumi Utiyama, Plant Gene Engineering Manual for Producing Transgenic Plants, Kodansha Scientific; page 71-74, ISBN4-06-153513-7 c3045. A tobacco leaf disc (approximately 0.5 g) from each transformed plant obtained in Examples 4, 5 and 6 was pulverized using a homogenizer in an Eppendorf tube to which was subsequently added 0.5 ml of 2×CTAB (2% acetyltrimethyl ammonium bromide, 1% polyvinylpyrrolidone (PVP)) pre-warmed to 65° C. and this mixture was incubated at 65° C. for 5 min. 0.5 ml of chloroform/isoamylalcohol (24:1) was then added and the sample was gently mixed for 5 min. The sample was separated by centrifugation at 12,000 rpm (10,000×g) for 10 minutes, 0.5 ml of isopropyl alcohol was added to the upper phase and the sample was mixed. After centrifugation at 12,000 rpm (10,000×g) for 15 min the precipitate obtained was dissolved in 200 μl TE (10 mM Tris-HCl pH 8.0, 1 mM EDTA). RNase was added so as to obtain a final concentration of 10 μg/ml and the sample incubated for 30 min at 37° C. to degrade RNAs. Following RNase treatment a mixture of equilibrated phenol/chloroform/isoamylalcohol (25:24:1) was added, the sample thoroughly mixed and the upper phase was collected. {fraction (1/10)}th volume of 3 M sodium acetate (pH 5.2) and 2.5 volumes of ethanol were added, the sample was thoroughly mixed and approximately 5 μg of genomic DNA was obtained by centrifugation at 12,000 rpm (10,000×g) for 5 min. 2. Confirmation of Gene Introduction by the PCR Method Using 50 μg of the genomic DNA obtained above as a template and synthetic DNA with the nucleotide sequences shown in SEQ ID NO:10 and SEQ ID NO:11 as primers, the promoter region was amplified using the PCR method (30 reaction cycles of 94° C. for 1 min, 55° C. for 2 min and 72° C. for 3 min). The PCR product obtained was analyzed by electrophoresis in 12% polyacrylamide gel (PAGE) (80 V, 1 hr.). By this procedure a DNA fragment corresponding to the promoter region (approx. 250 bp for Gbox10/−90/GUS and Gbox11/−90/GUS or 290 bp for −90/GUS) with the expected size was obtained. Example 8 Self-fertilization and Development of Genetically Homologous Line from Transformed Plant The transformed young plant generated in Example 4 was transferred to soil and cultivated to give rise to a transformed plant. At the time of anthesis (flowering) the plants were self-pollinated and seeds were obtained from the mature flowers. The seeds were sterilized for 5 min in 1% sodium hypochlorite, after which they were inoculated onto MS agar media containing 100 μg/ml kanamycin. Clones whose seeds all sprouted and developed after inoculation were selected. Example 9 GUS Gene Expression in Transformed Tobacco Plant Tissues GUS staining of the seeds, leaves and roots of the transformed tobacco plant obtained in Example 8 (containing either the plasmid, Gbox10/−90/GUS or Gbox11/−90/GUS or the plasmid −90/GUS or pBI121 as a control) was carried out according to the methods described in Hirofumi Uchiyama, Plant Gene Engineering Manual for Producing Transgenic Plant, Kodansha Scientific; page 68-70, 1990, ISBN4-06-153513-7 c3045 and Jefferson, Plant Mol. Biol. Rep. 5: 387-405 (1987). Seeds of the respective transformed tobacco plants were sterilized in 1% sodium hypochlorite and transferred to MS agar media containing 100 μg/ml kanamycin and allowed to develop seedlings for approximately two weeks. Activity staining was performed using 5-bromo-4-chloro-3-indole-β-D-glucuronic acid (X-Gluc) as a substrate and measuring the amount of the precipitated blue pigment (indigotin). Staining: Seeds from 40 individuals of each plant sample (transformed tobacco plants containing the plasmid, Gbox10/−90/GUS, Gvox11/−90/GUS, −90/GUS or pBI121 for comparison, or an untreated tobacco variety, SR-1 as a control) were sliced with a scalpel and immersed overnight at 37° C. in GUS staining solution (1 mM X-Gluc, 0.5 mM K 3 [Fe(CN) 6 ], 0.5 mM K 3 [Fe(CN) 6 ], 0.3% Triton X-100). Plant tissues were then transferred to ethanol, destained by several washes in ethanol and the amount of remaining blue pigment precipitate measured. The results of the GUS staining assay of seeds, young leaves and roots of transformed plants containing the plasmid of the present invention, Gbox10/−90/GUS, Gbox11/−90/GUS, −90/GUS or pBI121 for comparison, or the untreated tobacco variety SR-1 as a control are shown in Tables 1-3. TABLE 1 GUS staining in seeds Degree of staining (%) Plasmid A B C none SR-1 (control) 0 0  0 100 −90/GUS (Comparison 1) 0 0 75 25 Gbox10/−90/GUS (present plasmid) 83  0 17 0 Gbox11/−90/GUS (present plasmid) 3 59  31 7 pBI121 (35S/GUS) (Comparison 2) 30  50  10 10 A: the entire seed is darkly stained B: an area of dark staining emerging around the base (primordium) of the root in the seed C: the base of the root in the seed is stained none: no staining TABLE 2 GUS staining in leaves Degree of staining (%) Plasmid dark moderate light none SR-1 (control)  0 0 0 100 −90/GUS (Comparison 1)  0 0 0 100 Gbox10/−90/GUS (present plasmid) 100 0 0  0 Gbox11/−90/GUS (present plasmid)  70 24  3  3 pBI121 (35S/GUS) (Comparison 2)  67 11  11   11 “dark”, “moderate”, “light”, and “none” refer to the degree of staining TABLE 3 GUS staining in roots Degree of staining (%) Plasmid dark moderate light none SR-1 (control)  0 0 0 100  −90/GUS (Comparison 1)  0 4 22  74 Gbox10/−90/GUS (present plasmid) 100 0 0  0 Gbox11/−90/GUS (present plasmid)  53 27  10  10 pBI121 (35S/GUS) (Comparison 2)  75 0 0 25 “dark”, “moderate”, “light”, and “none” refer to the degree of staining In the non-transformed (untreated) tobacco variety SR-1 none of the tissues exhibited staining. In plants transformed with −90/GUS plasmid containing the minimal necessary elements for transcription initiation from the 35S promoter light staining was seen in the base (primordium) of the root. In plants transformed with the plasmid of the present invention Gbox10/−90/GUS and Gbox11/−90/GUS virtually all specimens exhibited heavy staining in the seeds, leaves and roots. Particularly very heavy and uniform staining was observed in the plant subjected to Gbox10/−90/GUS. Example 10 GUS Gene Expression in Transformed Carrot GUS staining was conducted with leaves of the transformed carrot plant (containing the plasmid of the present invention, Gbox10/−90/GUS and pBI121 as a control) in which the presence of introduced gene was confirmed in Example 7 according to the method as described in Example 9. No staining was observed for the leaves of untreated untransformed carrot. Heavy staining was observed for the leaves of most of the test specimen of the plant transformed with the plasmid of the present invention, Gbox10/−90/GUS. Example 11 GUS Gene Expression in Transformed Rice Expression of GUS gene in the leaves of the small transformed rice plant (height 10 cm) prepared in Example 6 or the gene was confirmed in Example 7 was examined according to the GUS staining method as described in Example 9. While little staining was observed in the leaves of untreated untransformed rice, strong GUS activity was observed in the leaves of the plant transformed with Gbox10/−90/GUS. The DNA length of the gene casette comprising a promoter functional in plant cells, a structural gene of interest and a terminator functional in plant cells contained in the present plasmids Gbox10/−90/GUS and Gbox11/−90/GUS is 2,300 bp (140 bp for the promoter), while that of pBI121 (35S/GUS) is 2,960 bp (800 bp for the promoter), and the gene casette to be introduced became about 22% shorter (83% for the promoter). The composition of the media used in the Examples are described below. 1. MS Agar Media 34.7 g of MURASHIGE AND SKOOG (Flow Laboratories) is dissolved in 1 liter of distilled water and the pH of the solution was adjusted to pH 5.8 with 1 M KOH. 8 g of agar are added and the mixture is sterilized by autoclaving. 2. MS-NB Agar Media 0.1 mg/ml of 1-naphthalene acetic acid (NAA) and 0.1 mg/ml 6-benzylaminopterin (BA) are added to MS agar media. 3. LS Media 34.7 g of MARUSHIGE AND SKOOG (Flow Laboratories) and 30 g of sucrose were dissolved in 1 liter of distilled water and the pH of the resulting solution was adjusted to pH 5.8 by 1 M KOH and sterilized by autoclaving. 4. L Agar Media Media obtained by adding agar 8 g/L to LS media. 5. LS-D Agar Media Media obtained by adding 1.0 mg/L of 2,4-dichlorophenoxyacetic acid to LS agar media. 6. LS-D2 Agar Media Media obtained by adding 2.0 mg/L of 2,4-dichlorophenoxyacetic acid to LS agar media. 7. L Broth Media 10 g of Bacto-tryptone (Difco), 5 g of Bacto yeast extract (Difco), and 10 g of NaCl are dissolved in 1 liter of distilled water, the pH is adjusted to 7.0 with 5 M NaOH and the media is sterilized by autoclaving. According to the present invention efficient expression of a gene of interest can be attained.

The present invention relates to a promoter functional in plant cells and plasmid that can regulate efficient expression of a gene of interest in plant cells. The promoter comprises a G-Box element, which enhances expression of an operably linked gene of interest in plants or plant cells. A further object of the invention is the plasmid pGbox10 or pGbox11, as well as plants or plant cells transformed with said plasmid.

REFERENCE TO RELATED APPLICATION This application is a continuation-in-part of application Ser. No. 784,154, filed Oct. 28, 1991. FIELD OF THE INVENTION The present invention relates to novel antimicrobial compositions and, more particularly, to a class of compounds having specific quaternized amine compounds linked to specific phosphate esters which exhibit broad spectrum antimicrobial activity and also virucidal activity referred to hereinafter as "antimicrobial/virucidal phospholipids". The phospholipids of the invention are well tolerated by human tissue making them suitable for use in the preparation of personal care, household cleaning, germicidal disinfectant and cleaning and like products which exhibit enhanced antimicrobial and virucidal characteristics. BACKGROUND OF THE INVENTION Phosphate ester and quaternary amine compounds are well known and have been widely used for many years for a variety of applications including those requiring surfactant properties. Known phosphate esters do not generally exhibit any antimicrobial characteristics, and while quaternary amine compounds are known in general to exhibit antimicrobial activity, such compounds are extremely irritating and thus have limited usefulness in personal care and cosmetic products. More recently, various betaine-type derivatives having, in general, quaternized alkyl amine groups and at least one phosphorus-containing anion in the molecule referred to hereinafter as "synthetic phospholipids", have been disclosed and suggested as, for example, in U.S. Pat. Nos. 4,215,064, 4,233,192 and 4,380,637 to Lindemann et al., U.S. Pat. Nos. 4,209,449, 4,336,385 and 4,503,002 to Mayhew et al., and U.S. Pat. Nos. 4,243,602, 4,283,542 and 4,336,386 to O'Lenick et al. These synthetic phospholipids are suggested as exhibiting an outstanding combination of surfactant characteristics as well as being well tolerated by human tissue, i.e., they exhibit exceptionally low ocular irritation and oral toxicity. While these known phospholipids have been found useful as surfactants in a variety of personal care, household cleaning and the like products, such products also require the incorporation of antimicrobial preservatives to inhibit microbial spoilage and increase shelf life, and there is no suggestion that any of these compounds exhibit virucidal activity. SUMMARY OF THE INVENTION In accordance with the present invention there has now been discovered novel phospholipid agents which surprisingly exhibit both excellent broad spectrum antibacterial and antifungal activity suitable for use as preservative and/or disinfectant agents in a variety of personal care compositions, household cleaning formulations and the like. These agents have also been found to possess potent virucidal activity making them particularly useful as a disinfectant, and for immobilizing and/or killing a variety of infectious viruses. The novel antimicrobial agents of the invention, which also exhibit virucidal properties, comprise particular synthetic phospholipid compounds that may be represented by the following general formula: ##STR3## wherein: x=1 to 3 or mixtures thereof; x+y=3; z=x a=0 to 2; B=O - or OM; A=Anion; M is a cation; R, R 1 and R2 are the same or different and are alkyl, substituted alkyl, alkyl aryl or alkenyl groups of up to 16 carbon atoms with the proviso that the total carbon atoms in R+R 1 +R 2 is between 10 and 24. It has been discovered that the particular synthetic antimicrobial phospholipids of the invention not only surprisingly and unexpectedly exhibit both broad spectrum bactericidal and fungicidal activity suitable for use as preservative and/or disinfectant agents in personal care and household products, but such phospholipid compositions surprisingly also exhibit potent virucidal activity making them useful, for example, as a disinfectant in hospitals and the like. Moreover, such agents are extremely well tolerated by human tissue, i.e., they exhibit exceptionally low ocular and skin irritation and oral toxicity, and can be used in product formulations containing nonionic, anionic, amphoteric and/or cationic components without significant inhibition or reduction of the required antimicrobial and/or virucidal activity. Thus, such agents may be formulated into a wide range of end products among which are germicidal cleaning compositions for hospitals and the like. The antimicrobial agents of the invention may also be used in combination with other known antimicrobial agents, when desired for particular applications, to enhance the antimicrobial and virucidal efficacy thereof. In another aspect of the invention, there is provided a method of inhibiting the growth of microorganisms in personal care, household cleaning and the like products which comprises incorporating in a personal care or household cleaning formulation an antimicrobially effective amount of an antimicrobial phospholipid compound of the general formula: ##STR4## wherein: x=1 to 3 or mixtures thereof; x+y=3; z=x; a=0 to 2; B=O - or OM; A=Anion; M is a cation; R, R 1 and R 2 are the same or different and are alkyl, substituted alkyl, alkyl aryl or alkenyl groups of up to 16 carbon atoms with the proviso that the total carbon atoms in R+R 1 +R 2 is between 10 and 24. In a still further aspect of the present invention, there is provided a personal care composition or a household cleaning composition which comprises a surface active agent and an antimicrobial effective amount of an antimicrobial phospholipid compound component of the general formula: ##STR5## wherein: x=1 to 3 or mixtures thereof; x+y=3; z=x; a=0 to 2; B=O - or OM; A=Anion; M is a cation; R, R 2 and R 2 are the same or different and are alkyl, substituted alkyl, alkyl aryl or alkenyl groups of up to 16 carbon atoms with the proviso that the total carbon atoms in R+R 1 +R 2 is between 10 and 24. In yet another aspect of the invention there are provided compositions for use in the killing and/or immobilizing a variety of infectious viral organisms including disinfectant protection which comprises a virucidally effective amount of a antimicrobial/virucidal phospholipid agent of the general formula: ##STR6## wherein: x=1 to 3 or mixtures thereof; x+y=3; z=x; a=0 to 2; B=0 - or OM; A=Anion; M is a cation; R, R 1 and R2 are the same or different and are alkyl, substituted alkyl, alkyl aryl or alkenyl groups of up to 16 carbon atoms with the proviso that the total carbon atoms in R+R 1 +R 2 is between 10 and 24; or a virucidal agent of the general formula: ##STR7## wherein: x is as hereinabove defined; x+y=3; z=x; a=0 to 2; B=0 - or OM; A is on Anion; M is a Cation; R 3 is an amidoamine moiety of the formula: ##STR8## wherein: R 7 is alkyl, alkenyl, alkoxy or hydroxyalkyl of from 5 to 21 carbon atoms each, or aryl or alkaryl of up to 20 carbon atoms; R 6 is hydrogen or alkyl, hydroxyalkyl or alkenyl of up to 6 carbon atoms each or cycloalkyl of up to 6 carbon atoms, preferably of from 2 to 5 carbon atoms, or polyoxyalkylene of up to 10 carbon atoms; R 4 and R 5 , which may be the same or different, are selected from alkyl, hydroxyalkyl, carboxyalkyl of up to 6 carbon atoms in each alkyl moiety, and polyoxyalkylene of up to 10 carbon atoms; in addition R 4 and R 5 taken together with the nitrogen to which they are attached may represent an N-heterocycle; and n is an integer from 2 to 6. As used herein the phrases "antimicrobial" and "inhibiting microbial growth" describes the killing of, as well as the inhibition or control of the growth of bacteria (gram positive and gram negative), fungi, yeasts and molds. As used herein the phrase "virucidal" describes the killing of as well as the immobilization of infectious virus organisms. DETAILED DESCRIPTION OF THE INVENTION The present invention is directed to novel phospholipid agents which surprisingly and unexpectedly exhibit excellent broad spectrum bactericidal and fungicidal activity and effectiveness, effectively inhibit the growth of a variety of bacteria, yeasts, and molds as well as possessing potent virucidal killing and/or immobilizing activity for a variety of infectious viruses. Moreover, such active agents may be used in combination with or in the presence of anionic, nonionic, amphoteric and/or cationic surfactants without inhibition of the antimicrobial and virucidal efficacy thereof and are virtually non-irritating to the skin and eyes; thus, such antimicrobial agents may be used in diverse formulations and applications. The novel antimicrobial/virucidal agents of the present invention comprise a class of synthetic phospholipid compounds which may be represented by the following general formula: ##STR9## wherein: x=1 to 3 or mixtures thereof; x+y=3; z=x; a=0 to 2; B=O - , OM; A=Anion; M is a cation; R, R 1 and R 2 are the same or different and are alkyl, substituted alkyl, alkyl aryl or alkenyl groups of up to 16 carbon atoms with the proviso that the total carbon atoms in R+R 1 +R 2 is between 10 and 24. The antimicrobial/virucidal phospholipid compounds described which, as indicated, exhibit broad spectrum antimicrobial as well as potent virucidal activity while being substantially non-irritating to humans, can be prepared by reaction of tertiary amines and phosphate esters corresponding to the amine and phosphate ester moieties in the above formula. Such compounds can be prepared by reacting the corresponding tertiary amine and phosphate ester reactants in the molar ratio of 1:1 to 3:1, and preferably from about 2.0:1 to 2.5:1, of amine to phosphate ester, for the time necessary for the amine to be completely reacted. Tertiary amines suitable for use in accordance with the practice of the invention can be represented by the general formula: ##STR10## wherein R, R 1 and R 2 are the same or different and are alkyl, substituted alkyl, alkyl aryl, or alkenyl groups of up to 16 carbon atoms with the proviso that the total carbon atoms in R+R 1 +R 2 is between 10 and 24. Exemplary tertiary amines include: tributylamine (di(hydroxyethyl)hexyl)-amine bis(2-hydroxyethyl)cocoamine N,N-dimethyl-dodecylamine N,N-dimethyl-tetradecylamine N,N-dimethyl-hexadecylamine N,N-dimethyl-cocoamine N,N-dimethyl-cetylamine dimethyl (C 8 -C 16 ) alkyl amine The phosphate ester reactants suitable for use in accordance with the practice of the invention can be represented by the general formula: ##STR11## wherein: x=1 to 3 or mixtures thereof; x+y=3; B=O - or OM; Hal=halogen. The phosphate ester intermediate may be prepared by known procedures wherein phosphoric acid and various phosphate salts, and preferably monosodium phosphate, are reacted in an aqueous medium with epichlorohydrin, generally in the molar ratio of from 1:1 to about 1:3, until the reaction is complete. As noted, the instant invention is based upon the discovery that the phospholipid compounds of the invention described above are effective in controlling the growth of bacteria, yeasts and molds in diverse formulations and applications such as cosmetic, toiletries, personal care, household and related products and materials. The phospholipid agents of the invention are not only effective antimicrobials for the destruction or control of fungi and bacteria that cause degradation and deterioration of diverse personal care and household product formulations, but also by their activity against the organisms that can reside and accumulate on various surfaces, they can provide utility in sanitizing, disinfecting and bacteriostatic applications. The antimicrobial activity of the compounds described above has been confirmed using standard laboratory techniques, including the Minimum Inhibitory Concentration (MIC) technique. They have been found effective, for example, in inhibiting bacteria including S. aureus, E. coli, P. aeruginosa and S. choleraesuis. They have also been found effective against yeast and mold including C. albicans and A. niger. In these tests it has been determined that the presence of anionic, nonionic, amphoteric and/or cationic materials did not inhibit the antimicrobial efficacy nor did a variety of inactivators commonly encountered in personal care and household applications. The broad spectrum preservative characteristics of the antimicrobial phospholipids of the invention in typical cosmetic formulations have also been established and confirmed. Specifically, molds and yeasts which may be inhibited include Aspergillus niger, Candida albicans plus various species of Penicillium, Tricholphyton, Alternaria, Gliocladium, Paecilomyces, Mucor, Fusarium, Geotrichum, Cladosporium and Trichoderma. Examples of the bacteria include Salmonella choleraesuis, Serratia marcescens, Klebsiella pneumoniae, Enterobacter aerogenes, Aerobacter aerogenes, Proteus vulgaris, Streptococcus faecalis, Pseudomonas aeruginosa, Escherichia coli, Staphylococcus aureus, Staphylococcus epidermidis, M. luteus, P. mirabilis, P. cepacia, P. stutzeri and A. hydrophilia. Another aspect of the present invention is the discovery that the antimicrobial phospholipid compounds surprisingly and unexpectedly exhibit significant spermicidal and antiviral activity which further enhances the utility of the compounds of the invention for a diversity of applications. The virucidal activity of the phospholipid compounds described above has been confirmed using test methodology according to U.S. Environmental Protection Agency guidelines for determining the virucidal efficacy of disinfectants intended for use on dry inanimate environmental surfaces (U.S. E.P.A. Pesticide Assessment Guideline, subdivision G, Product Performance, 198, Section 91-30 pp 72-76). Specifically, virucidal efficacy has been found against Human Influenza A virus; Herpes Simplex, type 2, virus; and the Human Immunodeficiency Virus (HIV). The phospholipid compounds described above have activity against bacteria, yeasts, molds as well as a variety of infectious viral organisms when employed at appropriate levels of concentration and may be used to inhibit growth or effectively destroy these organisms. It should be obvious that the required effective concentration or amount will vary with particular organisms and also on a number of other factors in particular applications. In general, however, effective antimicrobial response is obtained when the active agent is employed in concentrations ranging between five and 10,000 ppm (parts per million) and preferably between about 50 and 1,000 ppm. Generally, the concentration of the agent required for bactericidal activity will be lower than the concentration required for fungicidal activity and the concentration of the agent required for virucidal activity will generally be the same or higher than the concentration required for fungicida activity. For other applications, amounts of from 0.04% to about 5%, or higher, and preferably 0.07% to 3.0%, by weight of the active agent of the present invention is incorporated into a composition or sprayed onto or otherwise applied to a substrate to be treated in order to prevent growth of bacteria, yeasts and molds or killing and/or immobilizing infectious viral organisms. It will also be understood that the antimicrobial agents of the invention may be used in combination with other antimicrobial and/or virucidal materials. The compatibility of the phospholipid compounds of the invention with human tissue, i.e., dermal and eye tissue has also been tested. In these tests, 48 hour human patch dermal evaluations (5% in water), in vitro ocular evaluations (3% in water) and repeated insult patch tests (3% in water) determined that the compounds are substantially non-irritating to humans, they are safe and suitable for use in eye area products and are not a skin sensitizer to humans. While the phospholipid compounds hereinabove described exhibit broad spectrum antimicrobial as well as potent virucidal activity, certain other phospholipid compounds surprisingly have also been found to possess potent virucidal activity. Such compounds are compatible with anionic, nonionic, amphoteric and/or cationic materials without inhibition of their virucidal efficacy and exhibit low sensitivity to human tissue. Phospolipid compounds which are also suitable as a virucidal agent have the general formula: ##STR12## wherein: x is as hereinabove defined; x+y=3; z=x; a=0 to 2; B=0 - or OM; A is on Anion; M is a Cation; R 3 is an amidoamine moiety of the formula: ##STR13## wherein: R 7 is alkyl, alkenyl, alkoxy or hydroxyalkyl of from 5 to 21 carbon atoms each, or aryl or alkaryl of up to 20 carbon atoms; R 6 is hydrogen or alkyl, hydroxyalkyl or alkenyl of up to 6 carbon atoms each or cycloalkyl of up to 6 carbon atoms, preferably of from 2 to 5 carbon atoms, or polyoxyalkylene of up to 10 carbon atoms; R 4 and R 5 , which may be the same or different, are selected from alkyl, hydroxyalkyl, carboxyalkyl of up to 6 carbon atoms in each alkyl moiety, and polyoxyalkylene of up to 10 carbon atoms; in addition R4 and R5 taken together with the nitrogen to which they are attached may represent an N-heterocycle; and n is an integer from 2 to 6. The antimicrobial/virucidal phospholipid compounds of the invention may be incorporated in diverse personal care and household product formulations as, for example, a preservative therefore and/or as a disinfectant agent, and the incorporation of the compounds of the invention into such products can be done in accordance with standard practices. The active virucidal ingredients described can be diluted or otherwise mixed with solvents, dispersants, wetting agents, carriers and the like for topical or therapeutic use as a virucide in any desired application formulation such as liquids, sprays, etc. In connection with suitable modes of application for virucidal results, the phospholipid agents can be mixed with one or more pharmaceutically acceptable solid inert carriers. The invention will now be further illustrated by reference to certain specific examples which are provided herein for purposes of illustration only and are not intended to limit the scope therein. Example 1 925.6 grams of soft water are charged to a reaction vessel and heat is applied to 50° C. 554.4 grams of dimethyl cocoamine (C 12 --66%; C 14 --26%; C16--8%) are charged into the reaction vessel under good agitation and heat is applied to 90° C. An aqueous solution of 938.8 grams of 40% active 2-propanol, 1-chlorophosphate (3:1) are charged into the reaction vessel in four equal increments over 1.5 hours using good agitation while maintaining the temperature at 90°-95° C. Heating is continued at 90°-95° C. until the pH (10%) is 6.5 or less and the percentage of free tertiary amine is 0.5% maximum; approximately six to nine hours. The reaction mixture is then cooled to 80° C., 55.2 grams of 50% NaOH are charged into the reaction vessel and the reaction mixture is heated back to 90° C. Heating at 90 ° C. is continued until the percentage of NaCl is 6.9±0.2 %, approximately one hour. The reaction mixture is then cooled to 50° C. and the pH (10%) is adjusted to 7.0±0.5 with citric acid (approximately 9.7 grams). 22.1 grams of H 2 O 2 (35%) are charged to the reaction vessel with good agitation and heat is applied to 90° C. and maintained for one hour. The reaction mixture is then cooled to 50° C. and discharged. The product is a clear liquid having <0.5% free amine, a pH (10%) of 7.0±0.5 and a specific gravity @ 25° C. of 1.05. Example 2 682.4 grams of propylene glycol and 453.0 grams of water are charged to a reaction vessel and heat is applied to 50° C. 655.2 grams of dimethyl cetylamine are charged into the reaction vessel with good agitation and heat is applied to 90° C. An aqueous solution of 938.8 grams of 40% active 2-propanol, 1 chlorophosphate (3:1) are divided into four equal increments and charged into the reaction vessel over 1.5 hours while maintaining the temperature at 90°-95° C. Heating is continued at 90°-95° C. until the pH (10%) is 6.5 or less and the free tertiary amine is <0.5%, approximately six to nine hours. The reaction mixture is then cooled to 80° C. and 47.3 grams of 50% NaOH is added with good agitation. Heat is applied to 90° C. and maintained until the percentage of NaCl is 6.1±0.2%, approximately one hour. The reaction mixture is then cooled to 50° C. and the pH (10%) is adjusted to 7.0±0.5 with citric acid, approximately 4.7 grams being added. 25 grams of 35% H 2 O 2 are charged into the reaction vessel, heat is applied to 90° C. and maintained for one hour. The reaction mixture is then cooled to 50° C. and discharged. The product is a clear liquid having a specific gravity @25° C. of 1.05, a pH (10%) of 7.0±0.5 and Free amine of <0.5%. Example 3 The products of Example 1 and Example 2 are screened for antimicrobial activity using a modified Minimum Inhibitory Concentration (MIC) testing protocol. The initial screening is conducted using the following test organisms: S. aureus ATCC #6538 C. albicans ATCC #10259 A. niger ATCC #6275 Penicillium variable ATCC #XXXX The growth media used are Brain Heart Infusion Broth for bacteria and Sabouroud Broth for yeast and mold. A series of ten sequential two-fold dilutions of the test material is made in an appropriate growth promoting culture medium for each organism to be tested. A standard number of microorganisms are inoculated into each of the prepared dilutions containing the medium plus the test material. Inoculated tubes are incubated at appropriate temperature for 72 hours. Visual readings are taken after 24, 48 and 72 hours. The 72-hour incubated tubes are subcultured on agar media to verify inhibition of growth. Data are recorded as positive or negative for growth at each of the dilutions of the test material under evaluation. The minimum lethal concentration is defined as the smallest concentration of antimicrobial agent that, on subculture, either fails to show growth or results in a 99.9% decrease in the initial concentration of inoculum. Comparative MIC data of the initial screening test are reported in Table I. TABLE I______________________________________Test Organism Example I Sample Example II Sample______________________________________S. aureus 20 ppm 60 ppmC. Albicans 20 ppm 80 ppmA. niger 10 ppm 30 ppmP. variable 10 ppm 80 ppm______________________________________ An additional test panel is conducted to evaluate the products of Example 1 and Example 2. The further tests are conducted with Pseudomonas aeruginosa ATCC #15442, E. coli ATCC #8739 and Salmonella choleraesuis ATCC #10708. The MIC test protocol described above is used in conducting the additional test. Comparative MIC data of the additional screening test are reported in Table II. TABLE II______________________________________Test Organism Example I Example 2______________________________________P. aerugenosa 80 ppm 80 ppmE. coli 20 ppm 160 ppmS. choleraesuis 20 ppm 80 ppm______________________________________ As can be seen, both the Example 1 and Example 2 products exhibit significant antimicrobial properties. Example 4 A series of typical personal care products are prepared by standard practices using the following proportion of ingredients: ______________________________________Product A Shampoo______________________________________ Sodium Lauryl Sulfate 15.0% by weight Water 85.0% Antimicrobial Phospholipid variable (Example 1)______________________________________ Compositions are prepared with the following proportions of the product of Example 1. ______________________________________Test Sample Example 1 Product______________________________________A-1 0.00% by weightA-2 0.25% by weightA-3 0.50% by weightA-4 1.00% by weight______________________________________Product B Make-Up Foundation______________________________________a) Steareth - 20 1.5% by weight Pigment 15.0% by weight 0.5% Kelzan AR/1% NaCl 76.0% by weightb) Steareth - 2 2.5% by weight Isopropyl Myristate 2.0% by weight Hexyl Laurate 2.0% by weight Dow Fluid 200/100 cs 1.0% by weight Antimicrobial Phospholipid variable Pigment: White 13.50% by weight Red 0.15% by weight Brown 1.20% by weight Yellow 0.15% by weight______________________________________ Compositions are prepared with the following proportions of the product of Example 1. ______________________________________Test Sample Example 1 Product______________________________________B-1 0.00% by weightB-2 0.25% by weightB-3 0.50% by weightB-4 1.0% by weight______________________________________Product C Lotion______________________________________ a) Steareth - 20 2.0% by weight Water 87.5% by weight Product of Example 1 variable b) Steareth - 2 3.0% by weight Isopropyl Myristate 5.0% by weight Cetearyl Alcohol 2.5% by weight______________________________________ Compositions are prepared with the following proportions of the product of Example 1. ______________________________________Test Sample Example 1 Product______________________________________C-1 Product of 0.0% by weight Example 1C-2 Product of 0.1% by weight Example 1C-3 Product of 0.5% by weight Example 1______________________________________ Example 5 The personal care products of Example 4 are subject to Preservative Challenge Tests as follows: Aliquots of each test preparation are inoculated with separate representative mixed cultures of bacteria and fungi. Plate counts to determine survivors are performed at 0 time and after 3, 7, 14, 21 and 28 days of incubation. Bacterial samples showing a less than 10 recovery at 14 days are re-inoculated at days. Results are presented as surviving organisms present at each time interval per gram of product tested. PRODUCT A Inoculum a) Mixed bacteria: Pseud. Aeruginosa (ATCC 15442); E. coli (ATCC 8739 or 11229); S. aureus (ATCC 6536). Mixed fungi: A. niger (ATCC 9642); P. luteum (ATCC 8644); C. albicans (ATCC 10231). ______________________________________TESTSAMPLE DAYS BACTERIA FUNGI CONTROL______________________________________A-1 0 2,100,000 740,000 <10 3 17,500 4,750 <10 7 2,100,000 740,000 <10 14 2,100,000 740,000 <10 21* 2,100,000 740,000 <10 28 2,100,000 740,000 <10A-2 0 2,100,000 740,000 <10 3 24,200 1,900 <10 7 <10 <10 <10 14 <10 <10 <10 21* <10 <10 <10 28 <10 <10 <10A-3 0 2,100,000 740,000 <10 3 16,900 9,700 <10 7 <10 <10 <10 14 <10 <10 <10 21* <10 <10 <10 28 <10 <10 <10A-4 0 2,100,000 740,000 <10 3 23,700 1,620 <10 7 <10 <10 <10 14 <10 <10 <10 21* <10 <10 <10 28 <10 <10 <10______________________________________ *21-day Reinoculation NOTE: Control is an uninoculated sample for background count. Bacterial and Fungal Counts are as organisms recovered per gram of sample Test Day is the number of days after inoculation of the test sample. As can be seen, the antimicrobial product of Example #1 is highly effective against both bacterial and fungal challenges at a concentration of 0.25%. Moreover, the antimicrobial product of Example #1 is not adversely affected by anionics such as sodium lauryl sulfate. PRODUCT B Inoculum a) Mixed bacteria: Pseud. aeruginosa (ATCC 15442); E. coli (ATCC 8739 or 11229); S. aureus (ATCC 6536). b) Mixed fungi: A. niger (ATCC 9642); P. luteum (ATCC 9644); C. albicans (ATCC 10231). ______________________________________TESTSAMPLE DAYS BACTERIA FUNGI CONTROL______________________________________B-1 0 2,100,000 740,000 <10 3 2,100,000 740,000 <10 7 2,100,000 740,000 <10 14 2,100,000 740,000 <10 21* 2,100,000 740,000 <10 28 2,100,000 740,000 <10B-2 0 1,980,000 750,000 <10 3 57,000 4,200 <10 7 <10 120 <10 14 <10 1,420 <10 21* <10 5,300 <10 28 <10 7,400 <10B-3 0 2,100,000 740,000 <10 3 12,000 3,400 <10 7 <10 <10 <10 14 <10 <10 <10 21* <10 <10 <10 28 <10 <10 <10B-4 0 2,100,000 700.000 <10 3 3,000 <10 <10 7 <10 <10 <10 14 <10 <10 <10 21* <10 <10 <10 28 < 10 <10 <10______________________________________ *21-day Reinoculation NOTE: Control is an uninoculated sample for background count. Bacterial and Fungal Counts are as organisms recovered per gram of sample Test Day is the number of days after inoculation of the test sample. As can be seen, the antimicrobial product of Example #1 is highly effective against both bacterial and fungal challenges at a concentration of 0.50%. At 0.25%, the product of Example #1 is effective against the bacterial inoculum but failed to completely eradicate the fungi after initial reductions were noted. PRODUCT C Inoculum a) Mixed bacteria: Pseud. aeruginosa (ATCC 15442); E. coli (ATCC 8739 or 11229); S. aureus (ATCC 6536). b) Mixed fungi: A. niger (ATCC 9642); P. luteum (ATCC 9644); C. albicans (ATCC 10231). ______________________________________TEST CONTROLSAMPLE DAYS BACTERIA FUNGI (Uninoculated)______________________________________C-1 0 2,100,000 310,000 610 3 2,700,000 350,000 1,220 7 TNTC* TNTC TNTC 14 TNTC TNTC TNTC 21 TNTC TNTC TNTC 28 TNTC TNTC TNTCC-2 0 2,400,000 250,000 <10 3 <10 6,340 <10 7 <10 5,100 <10 14 <10 1,260 <10 21* <10 2,140 <10 28 <10 2,970 <10C-3 0 1,900,000 290,000 <10 3 <10 2,170 <10 7 <10 <10 <10 14 <10 <10 <10 21* <10 <10 <10 28 <10 <10 <10______________________________________ *TNTC Too numerous to Count *21day Reinoculation NOTE: Control is an uninoculated sample for background count. Bacterial and Fungal Counts are as organisms recovered per gram of sample Test Day is the number of days after inoculation of the test sample. As can be seen, Test sample C-3 (0.5% Product of Example #1) is found to effectively eliminate both bacterial and fungal challenges within seven days of inoculation. The product of Example #1 at 0.5% is capable of functioning effectively as a preservative as measured by the above test parameters. The antimicrobial test results clearly show the effectiveness of these products in preserving these systems. Noteworthy is the fact that product of Example #1 is not affected by anionics such as sodium lauryl sulfate. Example 6 The virucidal efficacy of the product of Example 1 against human influenza A virus is demonstrated in this example. In this test, virucidal efficacy of the test sample is evaluated by reduction in infectivity recoverable from a virus-contaminated surface after exposure to the use-dilution of the product. The test is conducted according to U.S. Environmental Protection Agency guidelines for determining the virucidal efficacy of disinfectants intended for use on dry inanimate surfaces (U.S.E.P.A. Pesticide Assessment Guidelines, Subdivision G: Product Performance, 1982, Section 91-30, pp. 72-76). In order for disinfectant efficacy to be claimed, the following criteria must be met in the test: 1. At least four logs of virus infectivity must be demonstrated, i.e. it must be possible to dilute the virus control four times 10-fold serially and still be able to detect infectious virus in the 10 -4 dilution. 2. The disinfectant must cause a 3 log reduction in virus titer. 3. There can be no detectable virus in the lowest non-toxic dilution of the virus-disinfectant sample. Human influenza A, strain A2/Hong Kong/8/68, ATCC VR-544, is the virus used in the study of this example. The virus suspension is prepared in allantoic fluid. The phospholipid compound used in this example is diluted for evaluation on the day of use 1:40 in sterile deionized water. Fertile chicken eggs incubated at 37 degrees C. are used which are candled on the day of inoculation; only live embryonated eggs being used. The embryonated eggs are inoculated after 10 days of incubation. The films of virus are made by placing 0.2 ml amounts of undiluted virus suspension on the bottoms of sterile glass Petri dishes and spreading. Films are held at room temperature (approx. 23 degrees C.) and ambient humidity, protected from direct light until dry (approximately 35 minutes). The dried virus films are treated with 2.0 ml of the use-dilution of the disinfectant sample for an exposure period of 10 minutes at approximately 23 degrees C. After exposure, the bottom of the dish is scraped with a rubber policeman to remove the virus disinfectant mixture. Concurrently with disinfectant treatment of one virus film, a parallel virus control film is resuspended in 1 ml of Phosphate-buffered saline (PBS). Assays for virus recovery are carried out by immediately making serial dilutions in PBS with the virus-disinfectant and virus control preparations and subsequently inoculating into embryonated eggs At least four (4) eggs are used per dilution. The eggs are inoculated with 0.2 ml volumes, and incubated at 37 degrees C. for approximately 72 hours with daily examination for mortality, and then cooled overnight at 4 to 6 degrees C. Allantoic fluids are collected from each egg and centrifuged for 10 minutes at approximately 800 xx g. Hemagglutination (HA) tests are carried out by mixing 0.5 ml of each fluid with 0.5 ml of 0.5% chicken erythrocytes (in PBS and observing for HA during the next one to two hours at room temperature. Cytotoxicity controls are run by diluting the use-dilution of the lot of disinfectant sample serially in PBS, and inoculating into embryonated eggs concurrently with virus-disinfectant mixtures. The viability of embryonated eggs is determined daily for three days of incubation at 37 degrees C. Viral and cytotoxicity titers are reported as -log 10 of the 50% titration endpoint for infectivity (ID 50 ) or toxicity (TD 50 ), as calculated by the method of Reed and Miuench (Amer. J. Hyg. 27: 493-497, 1938). Results of the study are reported in Table III. TABLE III______________________________________HUMAN INFLUENZA A VIRUSEvaluation of the PHOSPHOLIPID Sample for virucidalefficacy against dried virus after a 10-minute exposureto a 1:40 dilution in sterile deionized water. Hemagglutination (HA) Cytotoxicity (No. Positive/ ControlsDilution Inoculated) (No. Dead/No.Inoculated Control Sample + Virus Inoculated)______________________________________10.sup.-1 4/4 0/4 0/410.sup.-2 4/4 0/4 0/410.sup.-3 4/4 0/4 0/410.sup.-4 2/4 0/4 0/4Virus Titer 4/0 ≦0.5(-log.sub.10) ID.sub.50)HA AssayCytotoxicity Titer ≧0.5(-log.sub.10 TD.sub.50)-Reduction of virus ≧3.5titer by test sample(-log.sub.10 ID5.sub.0)-HA Assay______________________________________ Based on the results of infectivity and cytotoxicity assays shown in Table III, the Phospholipid example demonstrates virucidal activity against human influenza A. Infectivity is not detected in the virus-disinfectant mixture at the lowest nontoxic dilution. The reduction in virus titer for the phospholipid product of Example 1 is ≧3.5 log. Example 7 The virucidal efficacy of the product of Example 1 against Herpes Simplex, Type 2 is demonstrated in this example. The virucidal efficacy assay of this example generally employs the assay method of Example 6 except as noted. The virus employed is Herpes Simplex, type 2, ATCC VR-734 prepared in tissue culture medium. The cell cultures used are prepared from Vero cells obtained from Southern Research Institute with the cultures routinely grown in supplemented minimal essential medium (MEM). The cultures are grown and used as monolayers in disposable tissue culture labware at 37 degrees C in a humidified atmosphere of 5% CO 2 in air. After infection, cultures are held in maintenance medium containing the same ingredients with a 2% fetal calf serum. The reagents, disinfectant test solution and preparation of virus films are as described in Example 6. Treatment of Virus Films with Disinfectant: --Dried virus films are treated with 2.0 ml of the use-dilution of the disinfectant sample and allowed to remain in contact for a total exposure period of ten minutes at approximately 23 degrees C. After approximately the first 6.5 minutes of exposure, the bottom of the dish is scraped with a rubber policeman, and an aliquot of the virus-disinfectant mixture is immediately added to a Sephadex column for separation of virus from disinfectant by gel filtration. Concurrently with disinfectant treatment of one virus film, a parallel virus control film is resuspended in 2 ml of Phosphate buffered saline (PBS) and an aliquot is applied to a Sephadex column after 6.5 minutes. Sephadex gel filtration is performed generally by the method of Blackwell and Chen (J.AOAC 53: 1229-1236, 1970). The column filtrates are collected and diluted ten-fold serially for assay of infectivity. Assays for virus recovery are made using dilutions of each virus-disinfectant and control virus preparation. The dilutions are inoculated into cell cultures, at least four cultures per dilution being used. Cell monolayers are inoculated with 0.05 ml and incubated for one hour at 37 degrees C. After absorption, maintenance medium (0.2 ml) is added and cultures are incubated at 37 degrees C. Cultures are scored for cytopathic effects (CPE) at seven days after inoculation. Cytotoxicity controls of each batch of disinfectant sample are determined by placing 2.0 ml in the bottom of a sterile Petri dish containing a film of 0.2 ml PBS and after about 6.5 minutes an aliquot is filtered through Sephadex. The column filtrates are collected and diluted ten-fold serially for titration of cytotoxicity. Calculations of results are carried out as described in Example 6. The results of infectivity and cytotoxicity assays are reported in Table IV. TABLE IV______________________________________Cytopathic-CytotoxicEffects (No. Positive/No. Inoculated)Dilution CytotoxicityInoculated Control Sample + Virus Controls______________________________________10.sup.-1 4/4 0/4 0/410.sup.-2 4/4 0/4 0/410.sup.-3 4/4 0/4 0/410.sup.-4 2/4 0/4 0/4Virus Titer 4.0 ≦0.5(-log.sub.10) ID.sub.50)Cytotoxicity Titer ≧0.5(-log.sub.10 TD.sub.50)-Reduction of virus ≧3.5titer by test sample(-log.sub.10 ID.sub.5 0)-______________________________________ Example 8 In this example, the virucidal efficacy of the product of Example 1 is evaluated as measured by the reduction in infectivity of Human Immunodeficiency Virus, HTLF-III RF strain of HIV-1 using test protocols as described in Example 7. Preparation of the Starting Materials The RF Strain of HTLV-III human immunodeficiency virus (HIV) is used in this study. The Virus is produced by cultures of RF virus-infected H 9 cells (H9/RF) and is concentrated from supernatant culture fluid by high speed centrifugation by the following procedure: cells are first pelleted from a H9/RF culture by centrifugation at 600× g for 15 minutes at 4 degrees C. The supernatant culture fluid is transferred to 50 ml centrifuge tubes and centrifuged at 32,500× g. for 90 minutes at 4 degrees C. The supernatant is decanted and the virus pellet is resuspended in 1/100 the original volume of complete RPMI 1640 medium without fetal bovine serum. Resuspended virus pellets are kept at 4 degrees C. until used to prepare virus films. The disinfectant used in this example is diluted 1:40 on the day of use in sterile deionized water. Phosphate-buffered saline (PBS) is that of Dulbecco and Vogt, 1954. Films of virus are made by spreading 0.2 ml amounts of concentrated virus suspension over 28 cm 2 on the bottom of sterile glass Petri dishes. Films are held at room temperature (approx. 23 degrees C.) until visibly dry (approximately 45 minutes) and then incubated at 35-37 degrees C. in a dry oven for an additional 30 minutes to increase the level of dryness. Method of Determining Virucidal Efficacy of Disinfectant Treatment of Virus Films with Disinfectant: Dried virus films are treated with 2 ml of the diluted disinfectant and allowed to remain in contact for a total exposure period of 10 minutes at approximately 23 degrees C. After about 6.5 minutes of exposure, the treated virus films are filtered in a Sephadex column as described in Example 7. The column filtrates are diluted 10-fold for assay of infectivity. Treatment of Virus Control Films: A parallel virus film is resuspended in 2 ml of RPM 1640 medium without fetal bovine serum and antibiotics. After Sephadex filtration, the column filtrate is diluted 10-fold serially for assay of infectivity. Cytotoxicity Controls: The cytotoxicity of each batch of disinfectant test sample is prepared by placing 2 ml of the diluted disinfectant test sample in the bottom of a sterile Petri dish containing a film of dried PBS (0.2 ml). After about the first 6.5 minutes, an aliquot is filtered through Sephadex and subsequently diluted 10-fold serially for assay of cytotoxicity. Infectivity Assay: MT2 cells are indicator cells for infectivity assay. The MT2 cells are treated with polybrene (2 μg/ml) for 30 minutes at 37 degrees C., collected by centrifugation and plated in 96-well culture plates at approximately 1×10 4 cells per well in 0.15 ml of medium. Dilutions of each of the test and control groups are inoculated (0.05 ml/well) into four replicate cultures of MT2 cells and the cultures are scored for lytic cytopathic effects (CPE) after eight days of incubation at 37 degrees C. Viral and cyctotoxicity titers are expressed in this example as -log 10 of the 50 percent titration endpoint for infectivity (ID 50 ) or toxicity (TD 50 ), respectively, as calculated by the method of Reed and Muench. The results of infectivity and cytotoxicity assays are shown in Table V. TABLE V______________________________________CPE Assay with MT2 Cells (Day 8)Cytopathic-Cytotoxic Effects(No. Positive/No. Inoculated)Dilution CytotoxicityInoculate Control Sample + Virus Controls______________________________________10.sup.-1 Toxic 0/410.sup.-2 4/4 0/4 0/410.sup.-3 4/4 0/4 0/410.sup.-4 0/4 0/4 0/4Virus Titer 5.7 ≦1.5(-log.sub.10) ID.sub.50)Cytotoxicity Titer ≧0.5(-log.sub.10 TD.sub.50)-Reduction of virus ≦4.2titer by test sample(-log.sub.10 ID5.sub.0)-______________________________________ The results of infectivity and cytotoxicity demonstrated that the product of Example 1 possessed virucidal activity against HIV-1 in a CPE assay with MT2 cells. Example 9 The virucidal efficacy of various synthetic phospholipid compounds against human influenza A virus is demonstrated in this example. The synthetic phospholipid compounds evaluated in this example are: Product A--Cocamidopropyl PG--Dimonium Chloride Phosphate available commercially under the tradename PHOSPHOLIPID PTC from Mona Industries. Product B--Stearamidopropyl PG--Dimonium Chloride Phosphate available commercially under the tradename PHOSPHOLIPID SV from Mona Industries. In this Example, virucidal efficacy of Product A and Product B are evaluated by reduction in infectivity recoverable from a virus-contaminated surface after exposure to the use-dilution of the test products. The tests are conducted according to U.S. Environmental Protection Agency guidelines described in Example 6. Human influenza A virus, strain A/PR/834, ATCC VR-95 is used in the studies of this example. The virus suspension is prepared in tissue culture medium and is held in maintenance medium after infection containing the same ingredients in which the cultures are routinely grown but with 2% fetal calf serum instead of 10% serum. Virus films to be used are prepared as described in Example 6 as are the disinfectant product samples and phosphate-buffered saline (PBS) reagent. Treatment of virus films with disinfectant is carried out by treating dried virus films with 2.0 ml of the use-dilution of the disinfectant test samples and allowed to remain in contact for a total exposure period of 10 minutes at approximately 23 degrees C. After about the first 6.5 minutes of exposure, the bottom of the Petri dish is scraped with a rubber policeman, and an aliquot of the virus-disinfectant mixture is immediately added to a Sephadex column for separation of virus from disinfectant by gel filtration (see Example 7). Concurrently with disinfectant treatment of one virus film, a parallel virus control film is resuspended in 2 ml of PBS and an aliquot is applied to a Sephadex column after 6.5 minutes. The assays for virus recovery are carried out by making dilutions of each virus-disinfectant and control virus preparation and inoculating then into cell cultures. At least four cultures are used per dilution. Cell monolayers are inoculated with 0.05 ml and incubated for one hour at 37 degrees C. After absorption, maintenance medium (0.2 ml) is added and cultures are incubated at 37 degrees C. The cultures are scored for cytopathic effects (CPE) at seven days after inoculation. The cytotoxicity of each batch of disinfectant test sample is determined by placing 2.0 ml in the bottom of a sterile Petri dish containing a film of 0.2 ml PBS. After approximately 6.5 minutes, an aliquot is filtered through Sephadex. The column filtrates are collected and diluted 10-fold serially for titration of cytotoxicity. Viral and cytotoxicity titers are expressed as described in Example 7 and 8. The results of infectivity and cytotoxicity assays are shown in Table VI for both Product A and Product B. TABLE VI__________________________________________________________________________HUMAN INFLUENZA A VIRUSEvaluation of PRODUCT A AND PRODUCT B for virucidal efficacyagainst dried virus after a 10-minute exposureto a 1:40 dilution in sterile deionized water No. Dead/ (No. Positive/ No. Inoculated) No. Inoculated) Cytotoxicity Sample + Virus ControlsDilution Virus PRODUCT PRODUCT PRODUCT PRODUCTInoculated Control A B A B__________________________________________________________________________10.sup.-1 4/4 Toxic Toxic 4/4 4/410.sup.-2 4/4 Toxic Toxic 4/4 4/410.sup.-3 4/4 Toxic Toxic 4/4 4/410.sup.-4 2/4 0/4 0/4 0/4 0/4Virus Titer 5.7 ≦3.5 ≦3.5(-log.sub.10) TCID.sub.50)Cytotoxicity 3.5 3.5 3.5(-log.sub.10 TCTD.sub.50)Reduction of ≧2.2 ≧2.2virus titer bytest sample(-log.sub.10 TCID5.sub.0)__________________________________________________________________________ The results of infectivity and cytotoxicity demonstrate that Product A and Product B possess virucidal activity against human influenza A virus. Example 10 The virucidal efficacy of Product A and Product B of Example 9 against Herpes Simplex, Type 2 virus is demonstrated in this example. The procedure and ingredients of Example 7 are used in this study of the virucidal efficacy against Herpes Symplex Type 2, ATCC VR-734. The results of infectivity and cytotoxicity assays are shown in Table VII. TABLE VII__________________________________________________________________________HERPES SIMPLEX, TYPE 2Evaluation of PRODUCT A AND PRODUCT B for virucidal efficacyagainst dried virus after a 10-minute exposureto a 1:40 dilution in sterile deionized waterCytopathic-Cytotoxic Effects(No. Positive/No. Inoculated) Cytotoxicity Sample + Virus ControlsDilution Virus PRODUCT PRODUCT PRODUCT PRODUCTInoculated Control A B A B__________________________________________________________________________10.sup.-1 4/4 Toxic Toxic 4/4 4/410.sup.-2 4/4 Toxic Toxic 4/4 4/410.sup.-3 4/4 0/4 0/4 0/4 0/410.sup.-4 2/4 0/4 0/4 0/4 0/4Virus Titer 5.5 ≦2.5 ≦2.5(-log.sub.10) TCID.sub.50)Cytotoxicity 2.5 2.5(-log.sub.10 TCTD.sub.50)Reduction of ≧3.0 ≧3.0virus titer bytest sample(-log.sub.10 TCID5.sub.0)__________________________________________________________________________ Having now fully described the invention, it will be apparent to one of ordinary skill in the art that many changes and modifications can be made thereto without departing from the spirit or scope of invention as set forth herein.

There is provided a method for protecting substrates subject to contact by infections viral organisms by treating such substrates with virucidally effective amount of a composition containing a synthetic phospholipid of the formula: ##STR1## wherein: x is 1 to 3 or mixtures thereof; x+y=3; z=x; a=0 to 2; B=O - or OM; A is an Anion; M is a Cation; R 3 is an amidoamine moiety of the formula: ##STR2## wherein: n is an integer from about 2 to 6. R 4 , R 5 , R 6 , and R 7 are as defined in the specification.

BACKGROUND OF THE INVENTION 1. Field of the Invention The present invention relates to a method of manufacturing a semiconductor device, which in particular may have a high withstand voltage element and a low withstand voltage element mounted together on the same semiconductor substrate. 2. Background Information In recent years, liquid crystal displays have prevailed in the fields of personal computers and televisions, and their rapid growth in these fields is remarkable. Moreover, as liquid crystal displays have come to be used in cellular phones, digital cameras etc, it is expected that there will be even more demands for them in the future. A conventional liquid crystal panel requires high voltage in order to be operated. Therefore, a driver LSI for driving the liquid crystal panel needs high withstand voltage MOS (metal-oxide semiconductor) transistors. On the other hand, a logic circuit for digital processing needs an advanced logic process in order to obtain processing speed. Generally, in the logic processor, low withstand voltage MOS transistors are used. As opposed to this, in the driver LSI, elements having both high withstand voltage MOS transistors and low withstand voltage MOS transistors mounted on the same semiconductor substrate are used. Examples of general types of semiconductor devices having high withstand voltage MOS transistors and low withstand voltage MOS transistors mounted together on the same semiconductor substrate are exhibited in Japanese Laid-Open Patent Application No. 2000-150665 (hereinafter to be referred to as Patent Reference 1) and Japanese Laid-Open Patent Application No. 2000-200836 (hereinafter to be referred to as Patent Reference 2). FIG. 1A and FIG. 1B are diagrams showing structures of a general type of conventional semiconductor device 900 . FIG. 1A is a sectional view of the conventional semiconductor device 900 taken along a line I-I′, and FIG. 1B is an overhead diagram of the semiconductor device 900 . The I-I′ section of FIG. 1A is a section of the line I-I′ shown in FIG. 1B . Here, the same reference numbers are used for the same structural elements. As shown in FIG. 1A and FIG. 1B , the semiconductor device 900 has a high withstand voltage MOS transistor region 900 A and a low withstand voltage MOS transistor region 900 B. A MOS transistor formed in the high withstand voltage MOS transistor region 900 A (hereinafter to be referred to as high withstand voltage MOS transistor) has a gate oxide film 913 a and a gate electrode 914 formed on a silicon substrate 911 , sidewall spacers 916 formed on two sides of the gate electrode 914 , and a pair of source/drain regions 915 sandwiching a region underneath the gate electrode 914 in the silicon substrate 911 . On the other hand, like the high withstand voltage MOS transistor, a MOS transistor formed in the low withstand voltage MOS transistor region 900 B (hereinafter to be referred to as low withstand voltage MOS transistor) has a gate oxide film 913 b and a gate electrode 914 formed on the silicon substrate 911 , sidewall spacers 916 formed on two sides of the gate electrode 914 , and a pair of source/drain regions 915 sandwiching a region underneath the gate electrode 914 in the silicon substrate 911 . The high withstand voltage MOS transistor and the low withstand voltage MOS transistor are electrically separated from each other by a field oxide (a field oxide is also called an element isolating insulation film) 912 formed in the silicon substrate 911 . In the above structure, a boundary 900 a between the high withstand voltage MOS transistor region 900 A and the low withstand voltage MOS transistor region 900 B is positioned on the field oxide 912 . Now, with reference to FIG. 2A to FIG. 3B , a method of manufacturing the semiconductor device 900 according to prior art will be explained. FIG. 2A to FIG. 3B show manufacturing processes focusing attention on the section of the line I-I′ shown in FIG. 1B . First, as shown in FIG. 2A , field oxides 912 are formed in the p-type silicon substrate 911 using a well known STI (Shallow Trench Isolation) method for instance. By this arrangement, active regions and field regions are defined in the surface of the silicon substrate 911 . Next, by conducting a thermal oxidation treatment on the surface of the silicon substrate 911 , a gate oxide film 913 for the high withstand voltage MOS transistor is formed on the entire surface of the silicon substrate 911 as shown in FIG. 2B . Here, the gate oxide film 913 is normally formed to a thickness which is sufficient to not be damaged by an operating voltage. Generally, the gate oxide film 913 is formed to the thickness of around 30 to 50 nm (nanometer) for instance. Next, by conducting a known photolithographic process, a resist pattern R 901 is formed only in the high withstand voltage MOS transistor region 900 A. Then, the gate oxide film 913 in the low withstand voltage MOS transistor region 900 B is removed by a known etching method while using the resist pattern R 901 as a mask. By this process, the gate oxide film 913 A which is a part of the gate oxide film 913 remains only in the high withstand MOS transistor region 900 A as shown in FIG. 2C . The resist pattern R 901 remained on the gate oxide film 913 A is removed after the etching process is over. Next, by conducting a thermal oxidation treatment on the entire surface of the silicon substrate 911 , a gate oxide film 913 B for the low withstand voltage MOS transistor is formed in the low withstand voltage MOS transistor region 900 B as shown in FIG. 2D . Here, the gate oxide film 913 B is normally formed to a thickness which is decided depending on the operating voltage and performance expected from the low withstand voltage MOS transistor. Generally, the gate oxide film 913 B is formed to a thickness of around 2 to 7 nm for instance. Next, polysilicon is deposited on the entire surface of the silicon substrate 911 on which the gate oxide films 913 A and 913 B are formed, and processed by a known photolithographic method and an etching method to form the gate electrode 914 on the gate oxide film 913 A in the high withstand voltage MOS transistor region 900 A and the gate electrode 914 on the gate oxide film 913 B in the low withstand voltage MOS transistor region 900 B. Then, while using the gate electrodes 914 as masks, an etch back process is done on the entire surface of the silicon substrate 911 to remove the gate oxide films 913 A and 913 B except for the parts underneath the gate electrodes 914 . By these processes, a structure shown in FIG. 3A can be obtained. Next, an insulation film such as a silicon oxide film or a silicon nitride film is formed on the entire surface of the silicon substrate 911 using a known CVD (Chemical Vapor Deposition) method, after which an etch back process according to a known etching technique is performed on the insulation film to form the sidewall spacers 916 on the sides of the gate electrodes 914 respectively, as shown in FIG. 3B . Next, arsenic (As) ions are implanted into the silicon substrate 911 while using the field oxides 912 , gate electrodes 914 and the sidewall spacers 916 as masks, a pair of source/drain regions 915 are formed in the active region of each of the high withstand voltage MOS transistor region 900 A and the low withstand voltage MOS transistor region 900 B in a self-aligning manner, the pair of source/drain regions 915 being formed in a way sandwiching a region underneath the gate electrode 914 and the sidewall spacers 916 . Taking the processes described above, a semiconductor device having a low withstand voltage transistor and a high withstand voltage transistor formed on the same semiconductor substrate can be produced. However, in the above-described conventional manufacturing method, it has been noted as a problem that a step is produced in the upper part of the field oxide in the boundary between the high withstand voltage MOS transistor region and the low withstand voltage MOS transistor region. This is because in the etching process of FIG. 2C , over etching to the extent of about several dozen percent of the thickness of the gate oxide film 913 is done for the purpose of preventing variations in thickness to be made in the gate oxide film 913 B after etching process. Due to such over etching, the upper part of the field oxide 912 which is not covered by the resist pattern R 901 is also etched as shown in FIG. 2C . As a result, a step is formed in the upper part of the field oxide 912 in the boundary between the high withstand voltage MOS transistor region 900 A and the low withstand voltage MOS transistor region 900 B as can be seen in FIG. 2C . Normally, this step is about 50 to 100 nm high, although it depends on the thickness of the gate oxide film 913 . Such a step can be a cause of defective printing in the photolithographic process in forming the gate electrode 914 in the later process, and can be a cause of etching residuals of the polysilicon film. In addition to that, since the field oxide 912 becomes thinner, leakages, for instance, between transistors and wirings may be caused (hereinafter to be referred to as inter-field leakage). In the above described way, when there is a step in the upper part of the field oxide, problems such as open, short, leakage, etc. can occur, which leads to a problem in which normal operation of the semiconductor device becomes difficult. In view of the above, it will be apparent to those skilled in the art from this disclosure that there exists a need for an improved method of manufacturing a semiconductor device. This invention addresses this need in the art as well as other needs, which will become apparent to those skilled in the art from this disclosure. SUMMARY OF THE INVENTION It is therefore an object of the present invention to resolve the above-described problems and to provide a method of manufacturing a semiconductor device, which has a high withstand voltage element and a low withstand voltage element formed on the same semiconductor substrate and which does not form any step in a field oxide. In accordance with a first aspect of the present invention, a method of manufacturing a semiconductor device comprises the steps of: preparing a semiconductor substrate, the semiconductor substrate having first and second predetermined regions; forming a first field region surrounding the first predetermined region; forming a second field region surrounding the second predetermined region while a separating region exists between adjacent first and second field regions; forming a first insulation film on the semiconductor substrate; forming a resist pattern on the first insulation film, the resist pattern covering the first predetermined region and a part of the separating region; exposing the second predetermined region by etching the first insulation film using the resist pattern as a mask; forming a second insulation film on the second predetermined region; and forming gate electrodes on the first and second insulation films. In accordance with a second aspect of the present invention, a method of manufacturing a semiconductor device comprises the steps of: preparing a semiconductor substrate having first and second predetermined regions; forming a first element isolating insulation film encircling the first predetermined region; forming a second element isolating insulation film encircling the second predetermined region while a separating region exists between the first and second element isolating insulation films, the first and second element isolating insulation films being separated physically by the separating region; forming a first insulation film on the semiconductor substrate having the first and second element isolating insulation films; forming a protective film on the first insulation film; forming a resist pattern on the protective film, the resist pattern covering the protective film over the second predetermined region, and a part of the edge of the resist pattern being located over the separating region; exposing the first predetermined region and the first element isolating insulation film by etching the protective film and the first insulation film using the resist pattern as a mask; forming a second insulation film on the first predetermined region; exposing the second predetermined region and the second element isolating insulation film by etching the remaining protective film and the remaining first insulation film; forming a third insulation film on the second predetermined region; and forming gate electrodes on the second and third insulation films, respectively. In accordance with a third aspect of the present invention, a semiconductor device has a semiconductor substrate having first and second active regions, a first field region encircling the first active region, a second field region encircling the second active region, a separating region physically separating the first and second regions, gate insulation films formed on the first and second regions, and gate electrodes formed on the gate insulation films. These and other objects, features, aspects, and advantages of the present invention will become apparent to those skilled in the art from the following detailed description, which, taken in conjunction with the annexed drawings, discloses preferred embodiments of the present invention. BRIEF DESCRIPTION OF THE DRAWINGS Referring now to the attached drawings which form a part of this original disclosure: FIG. 1A and FIG. 1B are diagrams showing a structure of a semiconductor device 900 according to prior art; FIG. 2A to FIG. 2D are diagrams showing processes of forming the semiconductor device 900 according to a prior art manufacturing method; FIG. 3A and FIG. 3B are diagrams showing processes of forming the semiconductor device 900 according to the prior art manufacturing method; FIG. 4A and FIG. 4B are diagrams showing a structure of a semiconductor device 1 according to a first embodiment of the present invention; FIG. 5A to FIG. 5C are diagrams showing processes of forming the semiconductor device 1 according to a manufacturing method of the first embodiment of the present invention; FIG. 6A to FIG. 6C are diagrams showing processes of forming the semiconductor device 1 according to the manufacturing method of the first embodiment of the present invention; FIG. 7A and FIG. 7B are diagrams showing a structure of a semiconductor device 2 according to a second embodiment of the present invention; FIG. 8A to FIG. 8C are diagrams showing processes of forming the semiconductor device 2 according to a manufacturing method of the second embodiment of the present invention; FIG. 9A to FIG. 9C are diagrams showing processes of forming the semiconductor device 2 according to the manufacturing method of the second embodiment of the present invention; and FIG. 10A and FIG. 10B are diagrams showing processes of forming the semiconductor device 2 according to the manufacturing method of the second embodiment of the present invention. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS Selected embodiments of the present invention will now be explained with reference to the drawings. It will be apparent to those skilled in the art from this disclosure that the following descriptions of the embodiments of the present invention are provided for illustration only and not for the purpose of limiting the invention as defined by the appended claims and their equivalents. First Embodiment A first embodiment of the present invention will be described in detail with reference to the drawings. Structure FIG. 4A is a sectional view of a semiconductor device 1 according to the first embodiment of the present invention taken along a line II-II′, and FIG. 4B is an overhead diagram showing the semiconductor device 1 . The II-II′ section of FIG. 4A is a section of the line II-II′ shown in FIG. 4B . Here, the same reference numbers are used for the same structural elements. As shown in FIG. 4A and FIG. 4B , the semiconductor device 1 has a high withstand voltage MOS transistor region 1 A and a low withstand voltage MOS transistor region 1 B which are both semiconductor elements. An active region AR in the high withstand voltage MOS transistor region 1 A is defined by being electrically separated from the other regions by field oxides 12 A which are field regions FR. Likewise, an active region AR in the low withstand voltage MOS transistor region 1 B is defined by being electrically separated from the other regions by field oxides 12 b which are field regions FR. A high withstand voltage MOS transistor formed in the high withstand voltage MOS transistor region 1 A has a gate insulation film 13 a and a gate electrode 14 formed on a silicon substrate 11 , sidewall spacers 16 formed on two sides of the gate electrode 14 , and a pair of source/drain regions 15 sandwiching a region underneath the gate electrode 14 in the silicon substrate 11 . On the other hand, like the high withstand voltage MOS transistor, a low withstand voltage MOS transistor formed in the low withstand voltage MOS transistor region 1 B has a gate insulation film 13 b and a gate electrode 14 formed on the silicon substrate 11 , sidewall spacers 16 formed on two sides of the gate electrode 14 , and a pair of source/drain regions 15 sandwiching a region underneath the gate electrode 14 in the silicon substrate 11 . In the above structure, a p-type silicon substrate can be applied as the semiconductor substrate 11 for example. Furthermore, the field oxides 12 A and 12 b can be formed using an STI method for instance. However, the method of forming the field oxides 12 A and 12 b is not limited to the STI method, and can also be formed using a LOCOS (local oxidation of silicon) method for instance. The gate insulation film 13 a in the high withstand MOS transistor region 1 A can be an insulation film such as a silicon oxide film for instance. The gate insulation film 13 a should be formed to a thickness which is sufficient to not be damaged by an operating voltage, and such thickness may be around 30 to 50 nm for instance. On the other hand, the gate insulation film 13 b in the low withstand MOS transistor region 1 B can be an insulation film such as a silicon oxide film for instance, as with the gate insulation film 13 a . A thickness of the gate insulation film 13 b can be decided depending on the operating voltage and performance expected from the low withstand voltage MOS transistor, and it may be set to around 2 to 7 nm for instance. The gate insulation film 13 b is usually thinner than the gate insulation film 13 a. The gate electrodes 14 in the high withstand voltage MOS transistor region 1 A and the low withstand voltage MOS transistor region 1 B can be a polysilicon film including predetermined impurities, and they may be 200 to 300 nm thick for instance. The sidewall spacers 16 formed on the sides of each gate electrode 14 can be insulation films such as silicon nitride films for instance. However, it is preferable that the sidewall spacers 16 are made of a material which can be etched selectively under predetermined conditions with respect to the gate insulation films 13 a and 13 b , field oxides 12 A and 12 b , and the semiconductor substrate 11 . By choosing such material for the sidewall spacers 16 , it is possible to form the sidewall spacers 16 without having to form any resist patterns or the like for protecting the semiconductor substrate 11 , the sidewall spacers 16 , the field oxides 12 A and 12 b , the gate electrode 15 and so forth. For example, the sidewall spacers 16 can be formed without requiring any resist patterns or the like, under the conditions that the semiconductor substrate 11 is a silicon substrate, the field oxides 12 A and 12 b are silicon oxide films, the gate insulation films 13 a and 13 b are silicon oxide films, the sidewall spacers 16 are silicon nitride films, and a mixed gas of CHF 3 , Ar and O 2 with a mixture ratio of about 50:100:1 is used as an etching gas for processing the silicon nitride film formed on the semiconductor substrate 11 . In the active region AR of each of the high withstand voltage MOS transistor region 1 A and the low withstand voltage MOS transistor region 1 B, a pair of source/drain regions 15 are formed in the regions except the region underneath the gate electrode 14 and the sidewall spacers 16 , and the source/drain regions 15 are formed in a way which sandwich this region underneath the gate electrode 14 and the sidewall spacers 16 . In case of manufacturing a MOS transistor in which an n-type channel is formed, the source/drain regions 15 can be formed by implanting impurities having an n-type conductivity to the extent that the dose amount becomes around 2.0 to 5.0×10 12 /cm 2 . Here, arsenic (As) ions, for instance, can be used as the n-type impurities. On the other hand, in case of manufacturing a MOS transistor in which a p-type channel is formed, the source/drain regions 15 can be formed by implanting impurities having a p-type conductivity to the extent that the dose amount becomes around 2.0 to 5.0×10 12 /cm 2 . Here, boron (B) ions, for instance, can be used as the p-type impurities. In the above-described structure, as shown in FIG. 4A and FIG. 4B , the semiconductor device 1 has a semiconductor layer 1 b between the field oxide 12 A which defines the active region AR for the high withstand voltage MOS transistor region 1 A and the field oxide 12 b which defines the active region AR for the low withstand voltage MOS transistor region 1 B. This semiconductor layer 1 b is a region where the semiconductor substrate 11 is exposed, and it serves as a separating region which physically separates the field oxide 12 A and the field oxide 12 b . In this embodiment, a boundary 1 a between the high withstand voltage MOS transistor region 1 A and the low withstand voltage MOS transistor region 1 B is positioned on the semiconductor layer 1 b . Therefore, in this embodiment, in an exposure process in the manufacturing process of the semiconductor device 1 (which will be described later on) for instance, a photo mask with a layout enabling the boundary between the high withstand voltage MOS transistor region 1 A and the low withstand voltage MOS transistor region 1 B to be located on the semiconductor layer 1 b , and suitable exposure conditions for such purpose, are used. Here, the boundary to be located onto the semiconductor layer 1 b corresponds to the boundary 1 a. In this way, this embodiment realizes the structure in which the boundary 1 a between the high withstand voltage MOS transistor region 1 A and the low withstand voltage MOS transistor region 1 B can be located on the semiconductor layer 1 b , which is the exposed semiconductor substrate 11 , but not on the field oxide, i.e. the silicon oxide film. Therefore, in this embodiment, for instance, in a gate insulation film patterning process (q.v. FIG. 5C ) to be described later on, a boundary (which corresponds to 1 a ) between a region to be etched and a region not to be etched will not be disposed on the field oxide. As a result, it is possible to prevent any step from being formed in the upper part of the field oxide. Furthermore, as it will be mentioned later on, the semiconductor layer 1 b is a region where predetermined impurities are doped as with the active regions AR in the semiconductor substrate 11 . Therefore, by applying a predetermined electric potential to the semiconductor layer 1 b , possible inter-field leakage can be prevented. Accordingly, by having the structure according to this embodiment, it is possible to realize the semiconductor device 1 which is capable of preventing wire disconnection which can be caused by defective printing in the photolithographic process, or in other words, the semiconductor device 1 which is capable of preventing problems such as occurrences of open, short, leakage, etc. Manufacturing Method Now, a method of manufacturing the semiconductor device 1 according to the first embodiment of the present invention will be described in detail with reference to the drawings. FIG. 5A to FIG. 6C are diagrams showing processes of manufacturing the semiconductor device 1 according to the first embodiment of the present invention. With respect to FIG. 5A to FIG. 6C , each process will be described in terms of a section taken along the line II-II′ shown in FIG. 4B . First, as shown in FIG. 5A , field oxides 12 A and 12 B are formed in the upper parts of the semiconductor substrate 11 using a known STI method for instance. Thereby, the active region AR in the high withstand voltage MOS transistor region 1 A and the active region AR in the low withstand voltage MOS transistor region 1 B are defined, and the semiconductor layer 1 b which is an exposed semiconductor substrate 11 between the field oxides 12 A and 12 B (i.e. two field regions FR) is also defined. However, the method of forming the field oxides 12 A and 12 B is not limited to the STI method, but can also be a LOCOS method for instance. Next, by conducting a thermal oxidation treatment on the semiconductor substrate 11 , as shown in FIG. 5B , a gate insulation film 13 which has the same thickness as the gate insulation film 13 a for the high withstand voltage MOS transistor (e.g. around 30 to 50 nm) is formed on the whole upper surface of the semiconductor substrate 11 . Here, as for the conditions of the thermal oxidation treatment, for instance, the temperature is set at 850° C. and the heating time is set to around 30 to 40 minutes. Next, by conducting a known photolithographic process, a resist pattern R 1 is formed only in the high withstand voltage MOS transistor region 1 A. In this process, due to using a photo mask which enables the boundary 1 a shown in FIG. 4B to be focused onto the semiconductor layer 1 b , as shown in FIG. 5C , the side edge of the resist pattern R 1 is located on the semiconductor layer 1 b. Next, by using a known etching method, the gate insulation film 13 in the low withstand voltage MOS transistor region 1 B is removed. By this process, as shown in FIG. 5C , a gate insulation film 13 A which is a part of the gate insulation film 13 remains only in the high withstand voltage MOS transistor region 1 A. Here, in order to prevent the semiconductor substrate 11 from being damaged, it is preferable to use a wet etching method. In this wet etching process, for instance, the semiconductor substrate 11 having the gate insulation film 13 which is a silicon oxide film is soused in a hydrofluoric acid liquid of approximate 5% concentration for about 1 to 2 minutes. Moreover, in order to prevent etching residuals of the gate insulation film 13 from staying in the low withstand voltage MOS transistor region 1 B, it is preferable to conduct over etching to the extent of about several dozen percent of the thickness of the gate insulation film 13 . Here, an upper part of a field oxide in a region which is not covered by the resist pattern R 1 , i.e. the upper part of the field oxide 12 B in the low withstand voltage MOS transistor region 1 B, is also removed by this over etching. By this process, as shown in FIG. 5C , the field oxide 12 b of which upper parts is etched is formed. In addition, in the wet etching process using a hydrofluoric acid liquid as an etchant, it is possible to selectively etch the silicon oxide film (i.e. the gate insulation film 13 ) with respect to the silicon substrate (i.e. the semiconductor substrate 11 ). In this process, however, the surface of the semiconductor substrate 11 is etched slightly in over etching. Accordingly, a minute step 1 c is formed at the surface of the semiconductor substrate 11 as shown in FIG. 5C . After such etching process, the resist pattern R 1 on the remained gate insulation film 13 A is removed. Next, by conducting a thermal oxidation treatment on the semiconductor substrate 11 , a gate insulation film 13 B for the low withstand voltage MOS transistor is formed in the low withstand voltage MOS transistor region 1 B as shown in FIG. 6A . Here, the gate insulation film 13 B is formed to a thickness which is decided depending on an operating voltage and performance expected from the low withstand voltage MOS transistor. For instance, the gate insulation film 13 B is formed to the thickness of around 2 to 7 nm. As for the conditions of the thermal oxidation treatment, for instance, the temperature is set at 850° C. and the heating time is set to around 10 minutes. Next, by depositing polysilicon over the entire surface of the semiconductor substrate 11 on which the gate insulation films 13 A and 13 B are formed, a polysilicon film having a thickness of about 200 to 300 nm, for instance, is formed over the semiconductor substrate 11 . Then, by processing the polysilicon film using a known photolithographic method and a known etching method, gate electrodes 14 are formed on the gate insulation film 13 A in the high withstand voltage MOS transistor region 1 A and on the gate insulation film 13 B in the low withstand voltage MOS transistor region 1 B, respectively. Then, while using the gate electrodes 14 as masks, an etch back process is done on the entire surface of the silicon substrate 11 to remove the gate insulation films 13 A and 13 B except for the parts underneath the gate electrodes 14 . By these processes, a structure shown in FIG. 6B can be obtained. Next, an insulation film such as a silicon oxide film or a silicon nitride film is formed on the entire surface of the semiconductor substrate 11 using a known CVD method, after which an etch back process according to a known etching technique is performed on the insulation film to form sidewall spacers 16 on the sides of the gate electrode 14 respectively, as shown in FIG. 6C . Next, arsenic (As) ions are implanted into the semiconductor substrate 11 while using the field oxides 12 A and 12 b , the gate electrodes 14 and the sidewall spacers 16 as masks, by which a pair of source/drain regions 15 are formed in the active region of each of the high withstand voltage MOS transistor region 1 A and low withstand voltage MOS transistor region 1 B in a self-aligning manner, the pair of source/drain regions 15 being formed in a way which sandwich a region undernearth the gate electrode 14 and the sidewall spacers 16 , as shown in FIG. 4A . On the other hand, in order to form an electrode for controlling a substrate potential (which is also called a well potential), p-type impurities (e.g. boron (B) ions) are implanted into the semiconductor layer 1 b . By this arrangement, the electrical conductivity of the semiconductor layer 1 b can be improved. Taking the processes described above, a semiconductor device 1 having the low withstand voltage MOS transistor and the high withstand voltage MOS transistor formed on the same semiconductor substrate 11 can be produced. In the above-described way, according to this embodiment, the field oxide 12 A is formed in a way encircling the active region AR in the high withstand voltage MOS transistor region 1 A, the field oxide 12 b is formed in a way encircling the active region AR in the low withstand voltage MOS transistor region 1 B, and the semiconductor layer 1 b being the equivalent of the active region is formed between adjacent high withstand MOS transistor region 1 A and low withstand MOS transistor region 1 B. In this structure, since the boundary 1 a between the adjacent high withstand MOS transistor region 1 A and low withstand MOS transistor region 1 B is set on the semiconductor layer 1 b , it is possible to prevent any step from being formed in the field oxides 12 A and 12 b which electrically separate the high withstand MOS transistor region 1 A and the low withstand MOS transistor region 1 B. By such structure, in the photolithographic process and the like which form the gate electrodes 14 made of polysilicon for instance, it is possible to prevent problems such as defective printing, etching residual of the polysilicon film, etc. Moreover, since the semiconductor layer 1 b being the equivalent of the active region AR is laid out in a way encircling the high withstand voltage MOS transistor region 1 A (or the low withstand voltage MOS transistor region 1 B), by applying an arbitrary potential to this semiconductor layer 1 b , the possible occurrence of inter-field leakage can be prevented. Due to such effects, problems such as occurrences of open, short, leakage, etc. can be prevented in a semiconductor device in which a high withstand voltage MOS transistor and a low withstand voltage MOS transistor are formed in a single semiconductor substrate. Second Embodiment A second embodiment of the present invention will be described in detail with reference to the drawings. In the following, for the structures that are the same as the first embodiment, the same reference numbers will be used, and redundant explanations of those structure elements will be omitted. Structure FIG. 7A is a sectional view of a semiconductor device 2 according to the second embodiment of the present invention taken along a line III-III′, and FIG. 7B is an overhead diagram showing the semiconductor device 2 . The III-III′ section of FIG. 7A is a section of the line III-III′ shown in FIG. 7B . Here, the same reference numbers are used for the same structural elements. As shown in FIG. 7A and FIG. 7B , the semiconductor device 2 has a high withstand voltage MOS transistor region 2 A and a low withstand voltage MOS transistor region 2 B which are both semiconductor elements. An active region AR in the high withstand voltage MOS transistor region 2 A is defined by being electrically separated from the other regions by field oxides 22 a which are field regions FR. Likewise, an active region AR in the low withstand voltage MOS transistor region 2 B is defined by being electrically separated from the other regions by field oxides 22 b which are field regions FR. A high withstand voltage MOS transistor formed in the high withstand voltage MOS transistor region 2 A has the same structure as the high withstand voltage MOS transistor formed in the high withstand voltage MOS transistor region 1 A in the first embodiment, and in the high withstand voltage MOS transistor in the second embodiment, the gate insulation film 13 a is replaced with a gate insulation film 23 a . On the other hand, like the high withstand voltage MOS transistor, a low withstand voltage MOS transistor formed in the low withstand voltage MOS transistor region 2 B has the same structure as the low withstand voltage MOS transistor formed in the low withstand voltage MOS transistor region 1 B in the first embodiment, and in the low withstand voltage MOS transistor in the second embodiment, the gate insulation film 13 b is replaced with a gate insulation film 23 b. The gate insulation film 23 a is substantially the same as the gate insulation film 13 a in the first embodiment, but the formation process thereof differs from the formation process of the gate insulation film 13 a . Likewise, the gate insulation film 23 b is substantially the same as the gate insulation film 13 b in the first embodiment, but the formation process thereof differs from the formation process of the gate insulation film 13 b . These formation processes will be described later on in describing a manufacturing method, and therefore, detailed descriptions of those processes will be omitted in this past of the description. Furthermore, in the semiconductor device 2 of the second embodiment, the field oxide 12 A formed in the high withstand voltage MOS transistor region 1 A in the semiconductor device 1 of the first embodiment is replaced with the field oxide 22 a , and the field oxide 12 b formed in the low withstand voltage MOS transistor region 1 B in the semiconductor device 1 of the first embodiment is replaced with the field oxide 22 b. As can be seen from comparing FIG. 4A and FIG. 7A , unlike the field oxide 12 A, an upper part of the field oxide 22 a is removed and the amount removed in the field oxide 22 a is less than that in the field oxide 12 b . On the other hand, as can be seen from comparing FIG. 4A and FIG. 7A , the amount removed in the field oxide 22 b is less than that in the field oxide 12 b . These differences are provided by differences in the formation processes between the field oxides 12 A and 12 b and the field oxides 22 a and 22 b , respectively. These formation processes will be described later on as with the processes of the field oxides 22 a and 22 b. As described above, in the second embodiment, it is possible to reduce the removing amounts in the field oxide 22 a and 22 b . Furthermore, in this embodiment, the field oxides 22 a and 22 b are formed by etching field oxides 22 A and 22 B, which will be described later on, subject to the thickness of the silicon oxide film 27 , which will be described later on. Therefore, it is possible to conform the height of the upper face of the field oxide 22 a in the high withstand voltage MOS transistor region 2 A and the height of the upper face of the field oxide 22 b in the low withstand voltage MOS transistor region 2 B. Accordingly, in etching to remove a polysilicon film deposited over the field oxides 22 a and 22 b using a photolithographic method within the process of forming the gate electrode 14 , which will be described later on, it is possible to spread a margin to cope with displacements and defocuses (this margin is also called an exposure margin). Moreover, because the heights of the field oxides 22 a and 22 b are made even, it is possible to uniform the depths of contact holes formed over the gate electrodes 14 on the field oxides 22 a and 22 b . Accordingly, it is possible to spread a margin for etching conditions in forming the contact holes (this margin is also called an etching margin). Since the rest of the structure is the same as the first embodiment, a detailed description thereof will be omitted here. Manufacturing Method Now, a method of manufacturing the semiconductor device 2 according to the second embodiment of the present invention will be described in detail with reference to the drawings. FIG. 8A to FIG. 10B are diagrams showing processes of manufacturing the semiconductor device 2 according to the second embodiment of the present invention. With respect to FIG. 8A to 10B , each process will be described in terms of a section taken along the line III-III′ shown in FIG. 7B . First, as shown in FIG. 8A , field oxides 12 A and 12 B are formed in upper parts of the semiconductor substrate 11 using a known STI method for instance. Thereby, the active region AR in the high withstand voltage MOS transistor region 2 A and the active region AR in the low withstand voltage MOS transistor region 2 B are defined, and the semiconductor layer 1 b which is an exposed semiconductor substrate 11 between the field oxides 12 A and 12 B (i.e. two field regions FR) is also defined. However, the method of forming the field oxides 12 A and 12 B is not limited to the STI method, and can also be a LOCOS method for instance. Next, by conducting a thermal oxidation treatment on the semiconductor substrate 11 , a silicon oxide film 27 which is thinner than the gate insulation film 13 a for the high withstand voltage MOS transistor (e.g. around 10 to 20 nm) is formed on the whole upper surface of the semiconductor substrate 11 . Here, as for the conditions of the thermal oxidation treatment, for instance, the temperature is set at 850° C. and the heating time is set to around 20 minutes. Then, by depositing silicon nitride over the silicon oxide film 27 using a known CVD method for instance, a silicon nitride film 28 having a thickness of about 100 to 200 nm, for instance, is formed on the silicon oxide film 27 . By these processes, a structure shown in FIG. 8B can be obtained. The silicon nitride film 28 is a protective film with respect to a thermal oxidation treatment which will be described later on reference to with FIG. 9A . Therefore, any thickness of the silicon nitride film 28 is applicable as long it is a sufficient thickness with which the silicon oxide film 28 can protect the semiconductor substrate 11 from the thermal oxidation treatment. Next, by conducting a known photolithographic process, a resist pattern R 2 is formed only in the low withstand voltage MOS transistor region 2 B. In this process, as with the first embodiment, due to using a photo mask which enables the boundary 1 a shown in FIG. 7B to be focused onto the semiconductor layer 1 b , as shown in FIG. 8C , the side edge of the resist pattern R 2 is located on the semiconductor layer 1 b. Next, by using a known etching method, the silicon nitride film 28 and the silicon oxide film 27 in the high withstand voltage MOS transistor region 2 A is removed. By this process, as shown in FIG. 8C , a silicon nitride film 28 B which is a part of the silicon nitride film 28 and a silicon oxide film 27 B which is a part of the silicon oxide film 27 remain only in the low withstand voltage MOS transistor region 2 B. Here, in order to prevent the semiconductor substrate 11 from being damaged, it is preferable to use a wet etching method. In this wet etching process, for instance, the silicon nitride film 28 is etched by sousing the semiconductor substrate 11 having the silicon nitride film 28 in a thermal phosphoric acid liquid at a temperature of around 160° C. for about 30 to 40 minutes. The silicon oxide film 27 is etched by sousing the semiconductor substrate 11 having the silicon oxide film 27 in a hydrofluoric acid liquid of approximate around 5% concentration for about 1 to 2 minutes. Moreover, in order to prevent etching residuals of the silicon oxide film 27 from staying in the high withstand voltage MOS transistor region 2 A, it is preferable to conduct over etching to the extent of about several dozen percent of the thickness of the silicon oxide film 27 . Here, an upper part of a field oxide in a region which is not covered by the resist pattern R 2 , i.e. the upper part of the field oxide 12 A in the high withstand voltage MOS transistor region 2 A, is also removed by this over etching. By this process, as shown in FIG. 8C , the field oxide 22 a of which the upper part is etched is formed. However, since the thickness of the silicon oxide film 27 is thinner than the thickness of the gate insulation film 13 in the first embodiment (e.g. around 30 to 50 nm), and it may be around 10 to 20 nm for instance, the amount of the removed upper part of the field oxide 12 A is larger than the amount of the removed upper part of the field oxide 12 B which is removed in the etching process of the gate insulation film 13 in the first embodiment. That is, according to this embodiment, the amount of removed upper parts of field oxides can be reduced. In addition, in the wet etching using a thermal phosphoric acid liquid as an etchant, it is possible to selectively etch the silicon nitride film 28 with respect to the silicon oxide film 27 . Therefore, in this wet etching process, it is possible to ignore the thinning of the silicon nitride film 28 . Furthermore, in the wet etching process using a hydrofluoric acid liquid as an etchant, as with the first embodiment, it is possible to selectively etch the silicon oxide film 27 with respect to the silicon substrate (i.e. the semiconductor substrate 11 ). In this process, however, the surface of the semiconductor substrate 11 is etched slightly due to the over etching. Accordingly, a minute step 1 c is formed at the surface of the semiconductor substrate 11 as shown in FIG. 8C . After such etching processes, the resist pattern R 2 on the remained silicon oxide film 28 is removed. Next, by conducting a thermal oxidation treatment on the semiconductor substrate 11 , a gate insulation film 23 A for the high withstand voltage MOS transistor is formed in the high withstand voltage MOS transistor region 2 A as shown in FIG. 9A . Here, considering that the gate insulation film 23 A is to be thinned in the following process of etching the silicon oxide film 27 B, the gate insulation film 23 A should preferably be formed to a thickness that is thicker than a thickness that is sufficient to not be damaged by an operating voltage, by a portion of the gate insulation film 23 A to be thinned in the etching process of the silicon oxide film 27 B. The thickness sufficient to not be damaged by an operating voltage is around 30 to 50 nm for instance, and the portion of the gate insulation film 23 A to be thinned in the etching process of the silicon oxide film 27 B, in thickness, is around 11 to 22 nm for instance. Therefore, in this process, the gate insulation film 23 A should preferably be formed to a thickness of around 41 to 52 nm for instance. In addition, in this process, the silicon nitride film 28 B formed in the low withstand voltage MOS transistor region 2 B functions as a protective film with respect to a thermal oxidation treatment. Therefore, the gate insulation film 23 A should not be formed in the low withstand voltage MOS transistor region 2 B and the silicon oxide film 27 B should not be thickened. As for the conditions of the thermal oxidation treatment, for instance, the temperature is set at 850° C. and the heating time is set to around 30 to 40 minutes. Next, the silicon nitride film 28 B and silicon oxide film 27 B remaining in the low withstand voltage MOS transistor region 2 B are removed using the same wet etching method as in the process explained referring to FIG. 8C . Thereby, the semiconductor substrate 11 in the low withstand voltage MOS transistor region 2 B is exposed as shown in FIG. 9B . In this process, the gate insulation film 23 A is also etched. Thereby, a gate insulation film 23 C having a desired thickness (e.g. around 30 to 50 nm) is formed in the high withstand voltage MOS transistor region 2 A as shown in FIG. 9B . Furthermore, in this process, an upper part of the field oxide 12 B is removed by the over etching described above. Thereby, the field oxide 22 b in which the upper part is etched is formed. However, as described above, since the thickness of the silicon oxide film 27 is thinner than the thickness of the gate insulation film 13 in the first embodiment (e.g. around 30 to 50 nm), and it may be around 10 to 20 nm for instance, the amount of the removed upper part of the field oxide 12 B is larger than the amount of the removed upper part of the field oxide 12 B which is removed in the etching process of the gate insulation film 13 in the first embodiment. That is, according to this embodiment, the amount of the removed upper parts of the field oxides can be reduced. Moreover, according to this embodiment, it is possible to conform the removing thicknesses between the field oxides 12 A and 12 B. That is, it is possible to conform the height of the upper face of the field oxide 22 a in the high withstand voltage MOS transistor region 2 A and the height of the upper face of the field oxide 22 b in the low withstand voltage MOS transistor region 2 B. Next, by conducting a thermal oxidation treatment on the semiconductor substrate 11 , a gate insulation film 23 B for the low withstand voltage MOS transistor is formed in the low withstand voltage MOS transistor region 2 B as shown in FIG. 9A . Here, as with the first embodiment, the gate insulation film 23 B is formed to a thickness which is decided depending on the operating voltage and performance expected from the low withstand voltage MOS transistor. For instance, the gate insulation film 13 B is formed to a thickness of around 2 to 7 nm. As for the conditions of the thermal oxidation treatment, for instance, the temperature is set at 850° C. and the heating time is set to around 10 minutes. Next, by depositing polysilicon over the entire surface of the semiconductor substrate 11 on which the gate insulation films 23 C and 23 B are formed, a polysilicon film having a thickness of about 200 to 300 nm, for instance, is formed over the semiconductor substrate 11 . Then, by processing the polysilicon film using a known photolithographic method and a known etching method, gate electrodes 14 are formed on the gate insulation film 23 C in the high withstand voltage MOS transistor region 2 A and on the gate insulation film 23 B in the low withstand voltage MOS transistor region 2 B, respectively. Then, while using the gate electrodes 14 as masks, an etch back process is done on the entire surface of the silicon substrate 11 to remove the gate insulation films 23 C and 23 B except for the parts underneath the gate electrodes 14 . By these processes, a structure shown in FIG. 10A can be obtained. Next, an insulation film such as a silicon oxide film or a silicon nitride film is formed on the entire surface of the semiconductor substrate 11 using a known CVD method, after which an etch back process according to a known etching technique is performed on the insulation film to form sidewall spacers 16 on the sides of the gate electrode 14 respectively, as shown in FIG. 10B . Next, arsenic (As) ions are implanted into the semiconductor substrate 11 while using the field oxides 22 a and 22 b , the gate electrodes 14 and the sidewall spacers 16 as masks, by which a pair of source/drain regions 15 are formed in the active region of each of the high withstand voltage MOS transistor region 2 A and low withstand voltage MOS transistor region 2 B in a self-aligning manner, the pair of source/drain regions 15 being formed in a way which sandwich a region undernearth the gate electrode 14 and the sidewall spacers 16 , as shown in FIG. 7A . On the other hand, in order to form an electrode for controlling a substrate potential (which is also called a well potential), p-type impurities (e.g. boron (B) ions) are implanted into the semiconductor layer 1 b . By this arrangement, the electrical conductivity of the semiconductor layer 1 b can be improved. Taking the processes described above, a semiconductor device 2 having the low withstand voltage MOS transistor and the high withstand voltage MOS transistor formed on the same semiconductor substrate 11 can be produced. In the above-described way, according to this embodiment, the field oxide 22 a is formed in a way encircling the active region AR in the high withstand voltage MOS transistor region 2 A, the field oxide 22 b is formed in a way encircling the active region AR in the low withstand voltage MOS transistor region 2 B, and the semiconductor layer 1 b that is the equivalent of the active region is formed between adjacent high withstand MOS transistor region 2 A and low withstand MOS transistor region 2 B. In this structure, since the boundary 1 a between the adjacent high withstand MOS transistor region 2 A and low withstand MOS transistor region 2 B is set on the semiconductor layer 1 b , it is possible to prevent any step from being formed in the field oxides 22 a and 22 b which electrically separate the high withstand MOS transistor region 2 A and the low withstand MOS transistor region 2 B. By such structure, in the photolithographic process and the like in forming the gate electrodes 14 made of polysilicon for instance, it is possible to prevent problems such as defective printing, etching residual of the polysilicon film, etc. Moreover, since the semiconductor layer 1 b that is the equivalent of the active region AR is laid out in a way encircling the high withstand voltage MOS transistor region 2 A (or the low withstand voltage MOS transistor region 2 B), by applying an arbitrary potential to this semiconductor layer 1 b , the possible occurrence of inter-field leakage can be prevented. Due to such effects, problems such as the occurrence of open, short, leakage, etc. can be prevented in a semiconductor device in which a high withstand voltage MOS transistor and a low withstand voltage MOS transistor are formed in a single semiconductor substrate. Moreover, according to this embodiment, it is possible to reduce the removing amounts in the field oxide 22 a and 22 b . Furthermore, in this embodiment, the field oxides 22 a and 22 b are formed by etching field oxides 22 A and 22 B subject to the thickness of the silicon oxide film 27 . Therefore, it is possible to conform the height of the upper face of the field oxide 22 a in the high withstand voltage MOS transistor region 2 A and the height of the upper face of the field oxide 22 b in the low withstand voltage MOS transistor region 2 B. Accordingly, in etching to remove a polysilicon film deposited over the field oxides 22 a and 22 b using a photolithographic method within the process of forming the gate electrode 14 , it is possible to spread an exposure margin to cope with displacements and defocuses. Moreover, because the heights of the field oxides 22 a and 22 b are made even, it is possible to uniform the depths of contact holes formed over the gate electrodes 14 on the field oxides 22 a and 22 b . Accordingly, it is possible to spread an etching margin for etching conditions in forming the contact holes. Although the cases where an n-type high withstand voltage MOS transistor and an n-type low withstand voltage MOS transistor are mounted together on the same semiconductor substrate have been referred to in the above descriptions of the first and second embodiments, the present invention is not limited to this factor. For instance, by changing the impurities (ions) to be used, the present invention can be applied to a case where a p-type high withstand voltage MOS transistor and a p-type low withstand voltage MOS transistor are mounted together on the same semiconductor substrate, and a case where n-type and p-type high withstand voltage MOS transistors and n-type and p-type low withstand voltage MOS transistors are mounted together on the same semiconductor substrate. Although the case where two kinds of MOS transistors (i.e. a high withstand voltage MOS transistor and a low withstand voltage MOS transistor) are mounted together on the same semiconductor substrate has been referred to in the above descriptions of the first and second embodiments, i.e. the case where two kinds of gate insulation films with different thicknesses are formed on the same semiconductor substrate, the present invention is not limited to this factor. For instance, the present invention can be applied to a case where more than three kinds of gate insulation films with different thicknesses are formed on the same semiconductor substrate. Although the thermal oxidation treatment method is used for forming the gate insulation films ( 13 , 13 B, 23 A and 23 B) and silicon oxide film ( 27 ) in the above descriptions of the first and second embodiments, the present invention is not limited to this factor, and any method for forming a high resistance film having a desired thickness on the semiconductor substrate 11 can be applied as long as the formed film conforms to the spirit or the scope of the present invention. While the preferred embodiments of the invention have been described using specific terms, such description is for illustrative purposes only, and it is to be understood that changes and variations may be made without departing from the spirit or the scope of the following claims. This application claims priority to Japanese Patent Application No. 2005-69719. The entire disclosures of Japanese Patent Application No. 2005-69719 is hereby incorporated herein by reference. While only selected embodiments have been chosen to illustrate the present invention, it will be apparent to those skilled in the art from this disclosure that various changes and modifications can be made herein without departing from the scope of the invention as defined in the appended claims. Furthermore, the foregoing descriptions of the embodiments according to the present invention are provided for illustration only, and not for the purpose of limiting the invention as defined by the appended claims and their equivalents. Thus, the scope of the invention is not limited to the disclosed embodiments. The term “configured” as used herein to describe a component, section or portion of a device includes hardware and/or software that is constructed and/or programmed to carry out the desired function. Moreover, terms that are expressed as “means-plus function” in the claims should include any structure that can be utilized to carry out the function of that portion of the present invention. The terms of degree such as “substantially,” “about,” and “approximately” as used herein mean a reasonable amount of deviation of the modified term such that the end result is not significantly changed. For example, these terms can be construed as including a deviation of at least ±5% of the modified term if this deviation would not negate the meaning of the word it modifies.

A method of manufacturing a semiconductor device comprises the steps of: preparing a semiconductor substrate, the semiconductor substrate having first and second predetermined regions; forming a first field region surrounding the first predetermined region; forming a second field region surrounding the second predetermined region while a separating region exists between adjacent first and second field regions; forming a first insulation film on the semiconductor substrate; forming a resist pattern on the first insulation film, the resist pattern covering the first predetermined region and a part of the separating region; exposing the second predetermined region by etching the first insulation film using the resist pattern as a mask; forming a second insulation film on the second predetermined region; and forming gate electrodes on the first and second insulation films.

TECHNICAL FIELD [0001] The present invention relates to an interface for providing an input into a processor, i.e. a data processor, and can be used for numerous applications, including producing sound and music, providing signals representative of a wide variety of digital commands, providing the input of alphanumerical characters into a memory and/or a display, performing word or number processing functions or for generating signals providing remote movement or remote control of objects in real space or in a virtual environment, e.g. vehicle movement, or moving a remote robotic device or instrument, for example a surgical instrument. Especially, it relates to interfaces which simultaneously enable multiple kinds of inputs, including those denoting discrete commands, spatial location, and pressure-level. Numerous other applications will be evident to the skilled person on the basis of the discussion below. Background Field [0002] Currently, most user interfaces fall into one of two categories, which can be called “Discrete Control Interfaces (DCI)” which use a set of discrete switches which can register an on or off position to enable simple discrete inputs, and “Continuous Action Interfaces (CAI)”, which register spatial, pressure-based, or gestural movement in time to enable more complex inputs based on continuous movement. DCI include keyboards, keypads, and other interfaces that use direct analog (usually switch-based) controls that usually simulate a mechanical action, while CAI include touch screens, touchpads, other two-dimensional touch sensitive interfaces, and devices like a computer mouse, which use a rolling ball or some other continuous action apparatus that allows for continuous input. [0003] The advantages of DCI interfaces are (1) that they allow for clear discrete inputs, and (2) that they typically form a tactile input feedback system and thus do not rely on visual confirmation. In other words, they provide clear separate commands, and they give the user tactile information about which commands have registered, since the user can feel a responding pressure when he depresses a key, for example. These advantages relate not just to the kind of sensing device but also to the design of the input surface, the topmost part of the interface with which the user actually interacts. Here the springing quality of a typing keyboard allows the user to understand at the level of tactile perception that a key has been depressed, and the contours of the individual keys allows the user to make micro-adjustments to facilitate constant, fast, accurate typing without having to look at the keyboard. The various commands given by DCI can be memorised and can often be carried out without conscious thought and so information is usually faster and more accurately inputted than is the case with visually-based input control, especially when the DCI interface provides a tactile response. This can be seen by comparing the difference in speed between touch-typing (which depends on tactile-information feedback about finger location and activity, as well as habitual skill) and “hunting and pecking” typing (visually-based input control), or attempting to type on a touch screen interface. Another example is the difference between a virtuosic pianist (who uses tactile information about finger location and action, and habitual skill) vs. a beginner pianist who needs to look at the keys to orient himself as to which notes to play (visually-based tactile input control). These highly skilled activities depend on practice, and the reason practice is effective in these areas is that one can train one's muscular memory to repeat certain delimited tasks without conscious direction or control. In order for such training to be possible though, there are three requirements: a) the activity must not inherently require visual confirmation and direction (activities that require visual confirmation, like shooting a target, can of course also be practiced, but those involve a different form of practice involving hand/eye/body coordination); b) the physical interface must give positional tactile feedback (in the sense that a flat or merely decorated surface does not, and thus some kind of variation in surface, texture, or resiliency can consistently give the user something tactile to which to spatially orient his or her trained automatic muscular adjustments and correction); and c) these physical qualities of the interface have to be standardized and unchanging, so that they provide very similar tactile information in every instance. [0004] The disadvantage of DCIs is that they are limited in the types of input that can be made, especially when the goal is to input quantitative or continuous information, as opposed to qualitatively separate, distinct commands. [0005] On the other hand, CAIs have the advantage of allowing for continuous input and subtle or complex forms of information to be communicated very quickly. When, for example, a touch-sensitive interface is, or is connected to, a screen with changing information, it can be used to control a huge possible set of changing variables. The visual information can communicate innumerable options and then the pointer, whether one's finger or an arrow controlled by a mouse, can move to and select a particular set of movements and commands. Furthermore, this kind of input can simulate continuous human actions like handwriting or drawing in a way that would be impossible with a DCI. Their disadvantages however, are that (1) they are not as effective for processing a series of separate commands in rapid succession, in part because (2) they generally do not, in their current structure, provide a viable tactile feedback system in response to an input. [0006] It is therefore unsurprising that these two types of interface are often used in conjunction the computer keyboard and touchpad on most laptops is one example. The DCI aspect of the interface allows for discrete commands and rapid input for regular commands, like typing, while the CAI interface allows for more general control of the digital environment and any functions that require motion or continuity. Of course, while interfaces like a touchpad or a mouse are primarily CAI designed to register movement in time, they do allow for discrete inputs through their clicking or selecting function, and this combination of functions into one interface explains their immense usefulness. [0007] Multitouch gestures are well known, and these allow for other kinds of particular commands to be inputted into a CAI type interface, thus providing a further hybrid experience. A “gesture” signifies a single control of an interface, which may be a pattern simultaneously activated inputs such as the pressing of a letter key and the shift key to cause a capital letter to appear on a word processor screen. [0008] A certain degree of complexity is added in the case of pressure-sensitive interfaces. A variable pressure-sensing interface can be continuous with respect to pressure in the sense that it can register a continuous flow of data about different levels of pressure, while being discrete with respect the spatial distribution of pressure sensitive areas, depending on how many pressure sensors are used and how they are arrayed. [0009] Specific instances of relevant known equipment are: [0010] US 2003/0079549 discloses a pressure-sensitive device for accurately localising point sources of pressure and a method for determining the location force of a point pressure source on a pressure sensitive surface. It introduces the use of a internal membrane between the topmost surface in which the user inputs information, and the lower level in which the sensors actually read the data. This internal membrane can be made of a range of different materials and shapes, and can enable more accurate and balanced locational detection using fewer sensors than would otherwise be possible. It does not relate, however, to variations in the topmost surface other than protrusions to increase its hardness. Since the focus of the document is providing for more efficient ways of creating two-dimensional pressure sensing devices, (both in the sense that it gathers pressure data on an x and y axis, and in that the topmost surface is two-dimensional, it does not provide solutions for creating hybrid solutions for multiple pressure and spatially based discrete and continuous inputs. [0011] U.S. Pat. No. 6,703,552 discloses an electronic musical keyboard having a two-dimensional playing surface, i.e. it does not include discrete keys such as those present on pianos and similarly configured electronic musical keyboards. The keyboard continually measures each of the player's finger's pressure as it is pushed down on the keyboard to produce a sound. Because the surface is flat, it uses a program that automatically ‘guides’ the pitch to one of the main twelve tone notes if the player is close to one, which means that it has a somewhat inorganic sound, and it is very difficult to play it in a virtuosic way in this setting, because the computer is doing too much of the work. Also, it is difficult for a player to know precisely where his figures are on the keyboard because he has no information about the exact location of his fingers in relationship to the keys. It is a continuous, variable pressure and spatial-sensing interface that uses software algorithms to generate somewhat discrete effects. [0012] U.S. Pat. No. 7,394,010 discloses a musical instrument including an array of key switches for generating different tones. The switches are arranged in a two dimensional grid and sounds are produced by touching different key switches in the grid. [0013] EP 1100069 discloses a keyboard; the object of the keyboard is to simplify the electronic circuitry to detect the depression of keys. The keyboard includes a number of different buttons, each of which includes a number of contacts. The depression of each button closes a predetermined combination of the contacts and the combination differs from button to button. [0014] Therefore, by scanning the contacts periodically it is possible to tell which button is depressed. [0015] Electronic musical keyboards are, of course, very well known. They mimic a piano keyboard and include a number of different discrete notes. The notes can be modified to generate different sounds, e.g. a piano or an organ. However, the number of musical effects that is possible to achieve in real time with such a keyboard is limited. [0016] “Rollup” piano keyboards are also known. These take the form of a flexible substrate having a surface showing a two dimensional outline of a piano keyboard. Underneath each key, there is a contact which, when pressed, will produce a signal that instructs a processor and loud speaker to output a sound corresponding the key pressed. However, the operator gets very little tactile feedback and it is often difficult to know exactly where the figures are on the keyboard. [0017] Keyboard interfaces on touch screens are also well-known. They do not provide any tactile feedback that tells the user where his fingers are and so data entry will be done visually, i.e. not by touch typing. Also touch screens to not give a tactile feedback when the “key” is pressed sufficiently to register that the keystroke has resulted in the data being entered, although sometimes feedback is given in visual form (the key changes colour on the screen) or audible form (a “bleep” is sounded). DISCLOSURE OF THE INVENTION [0018] The present invention is defined in the accompanying claims. [0019] Broadly stated, the present invention provides an interface allowing a user to make an input into a processor. The interface includes a three dimensional input surface which provides tactile feedback to the user so that he knows where his fingers are on the interface. For ease of description, the input surface will sometimes be referred to in this specification as the “top” surface because it will sometimes be provided on the top of the interface (and other components of the interface will also be described by reference to such an interface orientation), but it should be remembered that the interface may be used in any desired orientation, e.g. with the input surface on a side or underneath surface, and the present specification covers the interface in any orientation. Furthermore, in many cases a given interface will consist of multiple surfaces where information could be inputted. Indeed, It could even be a totally enclosing surface with a top, bottom and sides. The input surface is simply the outermost part of the interface which is actually subject to the user's touch, in whatever form that takes. [0020] Soft resilient material, for example silicone rubber, is provided underneath (or within) the three-dimensional top input surface and indeed a surface of a layer of such soft resilient material could form the three-dimensional top surface. The top surface could be made, however, or any material that allowed for the diffusion of force or otherwise allowed for touch sensitivity. It could be made with a flexible OLED screen even, or a three-dimensional resistive or capacitive touch surface. [0021] The interface provides three distinct forms of tactile feedback to the user. Firstly, the texture, angle, and other characteristics of the three-dimensional top surface give the user immediate information about the location of the touch, in a way that would be impossible on a flat uniform surface where there is no tactile basis for spatial orientation. Secondly, the soft resilient material transmits forces back to the user to provide further tactile feedback to the user who will be able to sense the amount of pressure that he is applying to the interface. Because of the elasticity of the material, the response pressure increases exponentially with the distance of compression rather than simply pushing back with the same force as would a flat hard surface. This variation provides more usable tactile information to the user about slight changes in force input. Thirdly, the soft material amplifies the variation in the surface area of the tactile feedback. One of the ways we subjectively estimate the level of pressure with which we touch something, especially at very low levels of pressure input, is the variation in the amount by which the surface area of our finger, for example, touches a given surface. On a flat, hard surface, this difference is small, but with a softer material, the change is magnified, because as the muscle tissue in the finger compresses, so the does the material between the sensor and the input surface, and thus the surface area changes considerably. [0022] The downward force is diffused through the resilient material until it reaches the other side of the material, at which point it is applied to a sensor or an array of sensors, which may be one-dimensional or two dimensional. In some cases, the sensors could also be arrayed onto multiple underlying surfaces and thus constitute a set of two dimensional arrays (which could in turn be described as a three-dimensional array). If one sensor is provided, it must be capable of being able to discern what part of an input surface a user is manipulating and what forces/pressures are being applied at that location, i.e. it must have locational sensing capacity built into it. Each sensor is responsive to the touch of the user on the input surface, e.g. the sensors are pressure-sensitive, and they each output a signal in accordance with the touch-responsive parameter it measures, e.g. the force applied to it. [0023] The resilient material transmits the force applied to the top input surface to the sensors; because of the soft compliant nature of the material, the force is not only transmitted to the sensors located in line with the direction that the force is applied but also to other adjacent sensors. Thus, if a user applies a force to an area of the input surface of the interface, a wider area of the sensor array will register that force, causing a number of sensors to provide an output. The combination of distinct forces reading on a number of sensors will be referred to as a “signature”. Algorithms in the processor can respond to the make-up of the signature, including not only which sensors are “triggered”, i.e. provide a signal in response to applied pressure, but also the number of sensors, the forces they detect and the relative magnitudes of the outputs. For example, a strong force applied to the top surface will cause a relatively large number of sensors to be activated and produce output signals, as compared to a weaker force, whose effect is more narrowly transmitted by the soft resilient material to the sensors below, thereby triggering a smaller number of sensors. The simplest algorithms along these lines simply allow one to recover a particular location of touch more accurate than one could if one were touching the sensor array directly. If, for example, on a one dimensional array of sensors, three sensors in location A, B and C in a line give a reading of 100, 800, and 500, respectively one can weight the amount of force relative to the spacing of the sensors to recover an exact location of the input on the topmost surface. The total of the force readings is 1400, appropriately mapped or constrained depending on the application, and the location is given by (A*Aforce+B*Bforce+C*Cforce)/(total force), where Aforce, B*Bforce and Cforce are the readings from the sensors at locations A, B and C. If A, B and C, are spaced on a line at points 10 , 20 , and 30 , then the location of the input force is approximately 22.86. In other words, these kinds of algorithms can utilize a small number of sensors to give very accurate readings, especially when one allows for a significant amount of force diffusion. More complex algorithms can calculate similar data for two-dimensional arrays, and can, with sufficient resolution, also make calculations about different kind of shapes and gestures that a hand, for example, might make in interacting with a surface. This is particularly important with non-flat topmost surfaces because the algorithms have to translate data from surfaces on different angles and be able to distinguish between different kinds of touches. The particular algorithms needed for particular applications and circumstances would be apparent to a knowledgeable person in the field. These become more complex when set enacted over time, and with multiple inputs. In the previous example, if A continues to read at 100, B at 800, but C goes up to 600, that could represent an increase in the pressure and a slight movement towards C, but if at the same time sensors D and E also went up, for example to 800 and 100 respectively, the algorithm would have to be programmed to recognize that because the two highest force locations of B and D are separated by C (indicating two inputs centred on sensors B and D), C must be making a contribution to both B and D, and it can then either average that contribution, or it can hold the contribution of C at its previous level and contribute only the added amount of 100 to the roughly D centred input. This also indicates the limit of accuracy when applied to multiple inputs close together. At the same time, one can create gestural recognition in this context. For example, if A and E are initially the centre points of two inputs, and then B and D become the centre points that represents a pinching gesture, and the algorithms can be programs to register them as such. [0024] Although individual forces can be sensed as discussed above, changes in the signature of the sensors with time can also cause algorithms within the processor to generate certain effects, as will be discussed more broadly below. The processor processes the signals from the sensors and translates them into a processor output, which can be an electrical or optical signal, and the processor output can be used to produce a desired effect, e.g. a sound, a response on a screen, a recording of a piece of data in a memory, etc. Of course the processor output may be received by a different area of the processor, such as the recording of data within a memory area of the processor. It will be appreciated that a wide variety of other outputs and effects are possible. [0025] The three dimensional nature of the top surface not only provides tactile feedback to the user but different parts also provide different pressure signatures to the underlying sensors. Thus the forces sensed may depend on where on the input surface the user presses, e.g. pressing (a) on a peak of a raised area of the three dimensional input surface, (b) on a trough of the same surface or (c) on a shoulder or side of the peak will usually produce different sensor output signatures. This provides a product designer with many more options for producing different outputs (and different effects) based on the particular sensor signature detected. A given interface requires careful design and calibration regarding the nature of the surface, the thickness and hardness of the material, the size and sensitivity of the force sensors, and the nature of the algorithm of that translates a processor input into a processor output. [0026] The input surface of the interface may be an exposed surface of the soft resilient layer. However, the force transmission properties of the soft resilient layer may be modified by including one or more bodies either on the input surface or embedded within the soft resilient material. Such bodies may be harder or stiffer than the soft resilient material or indeed softer, e.g. air or gel pockets within the soft resilient material. They could also be different layers of similar material with different shapes, size and configurations, and levels of resilience. The softer or harder bodies will not only alter the sensor output signature but will also alter the tactile feedback sensed by the user. [0027] It is also possible to embed within the soft resilient material reactive devices or actuators that will initiate positive tactile feedback to the user based on output signals either from the sensors or from the processor. An example of such a device is one that will give an audible and/or a tactile response when a signal of predetermined characteristics is generated by the sensors. For example, if one or more sensors registers a force in excess of a threshold value, such a feedback could be produced. Silicone actuators which expand when a positive charge is run through them are known and can be arrayed with the three-dimensional resilient layer to create these haptic feedback effects. [0028] The sensors, in one embodiment, are a two-dimensional array of sensors and optionally are mounted on a hard rigid surface. [0029] As mentioned above, when a force is applied to the input surface over an area, this force is transmitted through the soft resilient material to the sensor array and the area of the array that will “feel” the pressing force will generally be larger than the area over which the user applies the force at the input surface. The degree of the spreading of the force will depend on several factors, including the softness and hardness of the resilient material, its thickness and also which part of the surface is pressed. As discussed above, it will also depend on whether or not other bodies are included within the soft resilient material. [0030] The minimum thickness of the soft resilient material may about 0.3 cm, e.g. at least 0.5 and generally at least 1 cm. The maximum may be 8 cm or even higher, e.g. up to 50 cm. A value that has been found to work well for most applications using finger pressure to provide the input to the interface is 2 to 4 cm. Furthermore, increasing the thickness of the soft resilient layer increases the bulk of the interface, which is itself a disadvantage. On the other hand, a thickness of less than 0.3 cm does not produce enough of a dissipation of the force applied to the input, so that the force applied to the sensor array will be very similar to that applied to the input surface. This reduces the range of pressure signatures that can be transmitted, which in turn reduces the range of effects that can be linked with those different signatures. [0031] As mentioned, the softness of the resilient material has an effect on the dissipation of the force from the input surface to the interface and it has been found that a very soft material having a shore hardness of 00-0001 to 10, e.g. 00-005 to 00-1, such as 00-01 to 00-1 provides satisfactory results for most applications using Finger pressure to provide the input to the interface. Using a precise shore hardness is important because if the material is too hard it creates too much resistance and latency and decreases the sensitivity of the sensors, and if it is too soft, it becomes too amorphous and doesn't provide sufficient tactile feedback. [0032] The sensors may be any sensors responsive to the user manipulating the input surface and may be pressure sensitive, in which case they may be piezo-electric crystals, strain gauges, for example, or may be made of a quantum tunnelling composite or they may be force sensitive resistors. Individual sensors and arrays of these types of sensors are widely commercially available; for example, the quantum tunnelling composite from Peratech Limited of Old Repeater Station, 851 Gatherley Road, Brompton on Swale, Richmond, North Yorkshire DL10 7JH United Kingdom and the force sensitive resistors from Interlink Electronics, Inc. of 546 Flynn Road, Camarillo, Calif. 93012, USA. [0033] One type of product where the present invention finds particular application is electronic musical instruments, partly because they produce a wide variety of different notes and but also because sound waves have a huge number of possible variables that create identifiably different sounds and an interface that can generate a large number of different signals is required to give rise to the variables. The most important identifiably different sounds are associated with rhythmic, and pitch, and volume variations. Rhythmic variations require the capacity for discrete inputs, while subtle pitch variation requires continuous input. Pitch is especially important—rhythmic variation can be provided by discrete input in time, i.e. by providing inputs that can be distinguished in time, while in the case of pitch, one needs to be able to input discrete distinct pitches, for example to play a scale in pitch, and also to create subtle pitch variations. At the same time, every note issues at a particular volume, and minor variations in volume create the basis for important differences in the emotive quality of the music. The problem is that when there is a need for both discrete and continuous pitch and volume variations, and a two-dimensional input surface cannot provide this—one has to choose between either inputting discrete pitches, which makes the continuous impossible, or continuous, which makes the discrete more or less impossible. [0034] Since the interface of the present invention enables seamless transitions for both discrete input (e.g. inputs to generate the notes of a chromatic scale) and continuous inputs (e.g. glissando and slide effects), it is ideally suited for the complexity of harmonic, dynamic and rhythmic variation. [0035] The interface may be in the form of a musical keyboard or other musical instruments played by touch. It may use the existing structure of an existing musical instrument and replicate them, i.e. the fundamental distribution of keys, strings, or buttons and their spatial intervals are the same as or similar to the original musical instrument, while also rendering it into a continuous 3D surface so that musicians can transfer their skills directly to a new instrument. [0036] Without a topmost surface that can be given a 3D form moulded to any shape and the diffusion of those forces onto a sensing array, it would be impossible to give a user the capacity to utilize existing muscle memory in service of a new, much more information rich interface. [0037] In one embodiment, the three dimensional input surface may have a wave-shape form where the peaks of the waves produce, when pressed, musical notes corresponding to the notes of a standard musical keyboard. In this way, the present invention can mimic a conventional keyboard in its operation. However, it has much greater versatility. For example, by pressing on one of the “peaks” or “crests” and vibrating a finger, an oscillating signature can be generated by the sensors, which will be interpreted by the processor as a vibrato. In addition, the shape of the surface means that a player can also play in the troughs, i.e. the areas between the crests, to produce microtonal pitches between any half or whole step. Since the input surface can be continuous, it is possible to produce smooth glissando effects on the keyboard. These particular effects can also he controlled through the software algorithms which can make the intervals in which one can play the 12 tone scale either wider, to avoid pitch bending, or narrower, to enable greater degrees of pitch bending. [0038] Another possibility is the provision of a further section of input surface that operates in the same way as the three dimensional input surface except that it is two dimensional; it may be provided either above or below the three dimensional input so that a user can easily transfer from one to the other or use them both simultaneously. It can be programmed to enable an operator to produce a smoother glissando effect (Portamento slider). A “palm effects slider” which allows the player to generate other sounds, or control programmable aspects of the timbre of the sound using her palm, in a way similar to certain hand-drums, like tablas is also possible. This variation between sections of the interface which are wavelike or corrugated, and sections which are flat and allow for long slides is an important component in maximizing the capacity for discrete and continuous inputs into the same interface. This allows for seamless transitions between the two. [0039] Other effects are also possible. For example, a peak of the three dimensional keyboard could be grasped or pinched, which will produce a particular signature of outputs from the sensors. This is especially the case if the grasped peak is manipulated in certain ways, i.e. pushed forward, sideways, up or down. The resulting signature can be interpreted by the processor to produce pre-set responses, for example to switch from outputting the music based on one musical instrument, for example a piano, to another form of musical instrument, for example, an organ, or even to jump between two samples of the same instrument, for example, legato and staccato samples of a violin. [0040] The interface can include algorithms that recognise specific input gestures corresponding to different kinds of musical samples in the same way that current keyboards use different samples depending on whether they have been hit harder or softer. In this case, however, a much wider set of gestures will allow the interface to produce a large number of simulations of a wide variety of instruments and effects. [0041] The principle of the present invention can be applied not only to musical keyboards but to simulation of any other musical instrument that are played by touch for example, a guitar or a violin. To simulate a violin or a guitar, it is possible to provide four or six peaks corresponding to the four or six strings of the instrument which will be responsive to a bowing, plucking or strumming action to produce a given sound. Similarly, different notes can be selected by the user touching the same peaks at a different part of the instrument, for example corresponding to the fret board on a standard stringed instrument. [0042] Other effects are possible, given the selectable nature of the material from which the interface is made. For example, if the interface is made in the shape of a guitar, the neck can be made flexible to produce, for example, a vibrato sound. [0043] The interface can be programmed in a variety of ways. One involves sending a MIDI signal each time a key is depressed, and using the initial pressure as the velocity level of the note, and the pitch bend function to control the exact pitch. If the signal then changes the MIDI note can remain on and the volume and pitch can be continuously adjusted. This is in effect more like volume and pitch ‘aftertouch’. Another way of programming the interface involves sending a new MIDI note with each pressure reading and simply changing the pitch level and the volume and velocity of every note. This method results in a more intuitive sound, but can give the sound a sense of texture. This texture can be corrected or accounted for by sending the MIDI notes at a frequency that approximates the natural waveform of the sound which one is trying to replicate. The interface can function with a variety of protocols and is not in any way limited to MIDI or these two approaches. [0044] Of course the present invention can also be used to create musical instruments with entirely new form factors, not merely for replicating the standard form of existing instruments in a new material. New forms can allow for new functions and new ways to integrate functions. [0045] The present invention is not limited to interfaces corresponding to musical instruments but can be applied to many other types of interfaces, for example computer keyboards to input alpha numeric characters or word or number processing functions. It is also possible to include within the interface a pointing device similar to a mouse. This may be provided for example, by a peak that can be grasped and pulled sideways, forward or back to control a pointer on a screen. [0046] The interfaces of the present invention can be used in other products which require a level of skill, since they can provide a high level of functionality that can only be achieved with a certain amount of practice. However, they could also be very useful for the elderly, or people with disabilities, because, using software, they can be tailored to the capacities of distinct users to respond to individual gestures. Their tactile quality makes them well suited to people who are either visually impaired or have more limited motor-control. [0047] One important aspect of the present invention is that it can be used as an interface for programmable gestures. In other words, a user could activate a setting on the processor the record a gesture and then touch the interface in a particular, unique way, and then instruct the processor that when that gesture is performed on the input surface the processor will activate a given command. This would allow individuals, and software programmers, a wide variety of options to play with, and a given interface could be tailored in software to precisely fit the needs of a given user. As mentioned above, this might have particular applications for the disabled, as well as for a broader market. [0048] The term “gestures” signifies a pattern of touching or pressing the input surface with the user's fingers or hand, including a pattern of moving the fingers or hand over the input surface. Thus gestures can include using the fingers to pinch, squeeze, push or pull part of the input surface in a particular way. The processor could be programmed to respond to a particular gesture by performing a particular command. Thus when any gesture is performed that resembles that stored in the processor, the processor could respond in a particular way. The gestures could be “programmed” into the processor by the user. A user could in essence “record” a given force signature as a particular “gesture” and then this could represent a particular kind of command. The interface would then be open to revision by the user not through a technical recoding of the fundamental program, but by simply saving a set of gestures into the processors memory that could tailor the functions of the interface to the particular needs and habits of a user. If for example, an upper-limb amputee wanted to use this kind of interface as a mouse, one kind of input into the top surface could control the moment of the arrow on screen, and another kind of “gesture” into the surface could be easily programmed to register a click. [0049] The three-dimensional quality of the topmost surface can also be moulded to fit exact irregular shapes and sizes. In the field of prosthesis, for example, there are currently robotic arms which allow for a high degree of articulated movement, including gripping at different levels of force, and making distinct gestures. Currently these prosthetic arms are usually controlled by electrodes. These send surface myographicelectric signals from the end points of the residual muscle tissue. These solutions often involve sequences of unintuitive movements and are never entirely natural and intuitive for the user. In the case of this invention, a soft material can be cast to form to the exact shape of the residual stump of the arm. Exact fitting is crucial for comfort and ease of use, and other pressure sensing solutions which do not involve complex three-dimensional input surfaces would not be suited for these kinds of applications. [0050] As mentioned above, the internal membrane between the three-dimensionally shaped input surface and the sensors can be made of a variety of different materials and can comprise a variety of different shapes and hardness to achieve a variety of desired effects. Using such an internal membrane to more accurately distribute pressure location is well known. Such techniques can be amplified and developed in the manner of the present invention, however, to detect small variations in exact locations on a three-dimensional input surface and to channel them to sensors lying on either a two or three dimensional substrate surface. [0051] This can be a kind of internal scaffolding structures which also enables or inhibits particular kinds of motion and input into the surface. Because the three dimensional internal flexible layer is in some application relatively thick, there are situations where one would want to restrict motion in a particular direction or strengthen a soft material. In such cases one can create a scaffolding or internal armature into a mesh. The exact form of the mesh can be varied according to the desired strength, restriction of motion, and sensing effects. Using an internal mesh within a softer material is a known technique, but in the present invention the mesh can also be used to simultaneously strengthen restrict the motion of a three dimensional shaped soft interface, and direct the forces on the three-dimensional input surface towards the sensor array. [0052] The scaffolding can pinpoint pressure sensitivity in particular locations and relay that directly to arrays of sensors thus giving very accurate readings of tiny variations on a three-dimensional input surface. Such a scaffolding system can in some instances also be cast with a resistive material like an ink-form of Quantum Tunnelling composite, and this can in turn make possible instances of invention in which there is no need for a rigid array of sensors underneath. The internal scaffolding can he made using rapid prototyping techniques which can allow for highly precise bespoke solutions tailored to meet individual needs. [0053] In the case or a prosthetic limb control system, for example, an internal scaffolding can be used to localize pressure output from particular muscle end-points onto particular sensors. This can enable to user to make intuitive gestures relating to the movement the missing hand, for example, and the program can then adjusted to recognize those gestures and send the right kind of data to the robotic arm. [0054] The interface of the present invention provides the responsive tactile feedback of keyboards and keypads—with the continuous input control of touch screen and touch pads. It then adds localized pressure sensitivity, particular three-dimensional forms to the input surface to allow for a range of desired inputs, and an internal scaffolding system where necessary to pinpoint force sensitivity to exactly where it is needed. The possibilities associated with bringing these categories of interface together into a new typology are enormous. BRIEF DESCRIPTION OF THE DRAWINGS [0055] The present invention will now be described in further detail, by way of example only, by reference to the accompanying drawings in which: [0056] FIG. 1 is a cross-sectional drawing of one possible formation that the components of a generalised interface in accordance with the present invention might take, integrating an input surface with a three-dimension shape, a flexible internal layer with a three-dimensional shape, and a two-dimensional sensor array. [0057] FIG. 2 is another cross-sectional drawing which provides a variation on FIG. 1 which highlights the fact that in some instances the sensors can also be array on three-dimensional shapes, (or multiple two dimensional surfaces which amount to a three dimensional array). [0058] FIG. 3 is a further cross-sectional drawings which provides another variation, and demonstrates that the input surface can be an entirely enclosing three dimensional shape and the sensors can be arrayed on a multi-surface three dimensional object. [0059] FIG. 4 is a cross-sectional drawing which shows an instance of the present invention in which the input surface is a recessed area and it is surrounded by the flexible force diffusing layer, and indeed a surrounding three-dimensional pressure sensing array. FIG. 5 is cross-sectional drawing which shows that the sensors can also be arrayed on complex three dimensional shapes, in this to provide another degree of freedom of movement—i.e. to create a force sensor that simultaneously responds to pushing and pulling. [0060] FIG. 6 is a cross-sectional drawing which shows the same arrangement as FIG. 5 , but adds in an internal scaffolding which is cast inside the resilient layer, the purpose of which is both increase the distribution of forces from particular points onto multiple sensors to give more accurate and detailed readings, and to make a more robust framework for manipulation in multiple directions. [0061] FIG. 7 is a schematic exploded drawing of one possible iteration of the whole system of the present invention. [0062] FIG. 8 is a view of one possible configuration of the input surface of a musical keyboard in accordance with the present invention; [0063] FIG. 9 is a close-up view of part of the surface in FIG. 8 . [0064] FIG. 10 is a flow diagram showing the general operational steps of the present invention of which the keyboard of FIG. 8 is an example; and [0065] FIG. 11 shows a keyboard-like tactile interface that can be used to input alphanumeric characters and/or to word or number process them, as well as control the cursor tracking, although other computer-related functions can also be driven from the interface. DETAILED DESCRIPTION [0066] Referring initially to FIG. 1 , which is a cross-sectional drawing of one possible formation that the components of a generalised interface in accordance with the present invention might take. It includes an illustrated point of contact 1 between a user and three-dimensionally shaped input surface(s) 2 . Between the three-dimensionally shaped input surface(s) 2 and two-dimensionally arrayed sensor array 4 there is a three-dimensionally shaped internal flexible layer 3 which Matches the contours of the input surface 2 and the sensor array 4 . The sensor array 4 is situated on hard backing surface(s) 5 . Also shown are illustrated lines of forces diffusion 6 from illustrated point of contact 1 and sensor array 4 . [0067] Referring now to FIG. 2 , it can be seen that the sensor array 4 can be situated on multiple two-dimensional surfaces which amount to a three dimensional array. The hard backing surface(s) 5 is also shaped to find the form of the sensor array, and as a result, the resilient layer 3 takes on a more complex three-dimensional character. The spatial relationship between the three-dimensionally input surface(s) 2 and the three-dimensional sensor array 4 can be optimized according to specific applications to maximize the range and accuracy of different kinds of force signature inputs and user gestures. [0068] FIG. 3 is another cross-sectional drawing which shows another possible arrangement of the components of the present invention. In this drawing, the three-dimensionally shaped surface(s) 2 entirely enclose the other components, including the three-dimensionally shaped flexible internal layer 3 and the three-dimensional sensor array 4 , and the hard backing surface(s) 5 . This arrangement allows for a range of different forms of user manipulation, including squeezing and stretching inputs into three-dimensionally shaped input surface(s) 2 , and also pressing the entire interface onto an external hard surface like a table. The three-dimensionally shaped surface(s) 2 can of course be one continuous shape or several connecting surfaces, and can have further textural details like ridges, protrusions, recesses, raised areas and any other complex shapes as is called for by the design or function of the particular application. [0069] FIG. 4 shows a cross-sectional view of a version of an interface of the present invention in which three-dimensional shaped input surface(s) 2 exist in an inverted configuration relative to the other drawings, such that three-dimensional shaped input surface(s) 2 are within three-dimensionally shaped flexible ‘internal’ layer 3 , and three-dimensionally shaped flexible ‘internal’ layer 3 is in turn within both three-dimensional sensor array 4 , and three-dimensionally arrayed hard backing surface(s) 5 . This arrangement can increase the surface area of the point of contact between the user and three-dimensional shaped input surface(s) 2 , which can allow for even more subtle registration of movement. In this arrangement it is sometimes necessary for three-dimensionally shaped flexible ‘internal’ layer 3 to be made of a softer, more gel-like material than in other arrangements, since there is less scope for the material to flex, and compress. [0070] Referring now to FIG. 5 , it can be seen that three-dimensional sensor array 4 can also be constructed with inverted three-dimensional shapes. This carries with it an important added benefit, which is that the it is one way of using the present invention to make a six degrees of freedom pressure-sensing device. One can pull upwards on three-dimensional shaped input surface(s) 2 and the force is translated through three-dimensionally shaped flexible ‘internal’ layer 3 onto the sensors of three-dimensional sensor array 4 , which are arrayed against the now complexly shaped hard backing surface(s) 5 . [0071] In FIG. 6 , the cross-sectional drawing shows an arrangement similar to that of FIG. 5 , with the addition of internal scaffolding system 7 . The internal scaffolding can be cast inside three-dimensionally shaped flexible ‘internal’ layer 3 , and can both increase the distribution of Forces from particular points onto multiple sensors to give more accurate and detailed readings, and make a more robust framework for manipulation in multiple directions. It can ensure that the pulling forces, for example, are appropriately translated onto the appropriate inverted sensors of three-dimensional sensor array 4 , and that the force diffusing qualities that come into play when pushing on the material don't reduce the accuracy in the case of pulling. The internal scaffolding system must have a certain degree of flexibility to function in this arrangement, but it must also be harder/stiffer than three-dimensionally shaped flexible ‘internal’ layer 3 . Nylon or higher shore hardness types of silicone are two materials which have been found to work well in these arrangements. In certain instances, force sensors and other resistive material can be built onto, or encapsulated within, this scaffolding, especially in cases where one desires a more completely flexible information-rich pressure sensitive interface. [0072] FIG. 7 , is a schematic exploded drawing of the components of a generalised interface in accordance with the present invention. It includes a three-dimensionally shaped flexible layer 13 of a soft resilient material, silicone rubber, having a three-dimensionally shaped input surface 12 . The body has a planar bottom surface 19 , which rests on an array of sensors 14 that is supported on a rigid surface 15 . Each of the sensors of the array 14 is pressure sensitive and can produce an output in accordance with the pressure exerted on it. The output from the various sensors is conducted via a lead to a microprocessor 16 . The microprocessor includes algorithms that respond to certain combinations of signals from the sensors of the array 14 to produce an output driving a component, which in this case is a loud speaker 17 . In many particular arrangements a small microprocessor, which could be, for example, an Arduino processor, will first interpret the data from the force sensors into a form that can be read as an input by a more advanced processor like a computer. The computer can then run the code that translates the basic numerical data of particular forces on particular sensors into a meaningful output. Such translation can be achieved by a variety of software solutions, again depending on the nature of the application and the desired output. For example, the data could be interpreted into particular outputs in a Java-based language like Processing, or a C++ based language like Open Frameworks, among many others. [0073] Referring now to FIG. 8 , which shows a simplified and schematic small section of three-dimensional shaped input surface(s) 12 it can be seen that the pressing 11 of a peak on the undulating surface 12 of the soft resilient body 13 transmits the force through the body to the planar surface 19 , which is in contact with an array of force sensors as seen in FIG. 7 . The force spreads out, as shown schematically by arrows, so that the area of the planar surface that “feels” the force on the three-dimensional shaped input surface(s) 12 has a greater area than the area pressed on the top surface. The microprocessor 16 ( FIG. 7 ) detects which sensors in the array are providing the signal, and the software algorithms can interpret the relative forces on each sensors to reconstruct the exact location of the input. Once the input data has reconstructed, the software then creates a set of software ‘objects’ which can keep track of a set of simultaneous inputs, and can compare each continuing input with the inputs from the previous loop of the program. Depending on the nature of the application and the desired level of discreteness or continuity between inputs, algorithms can be set to measure the level of continuity which is necessary for the program to interpret a given input as being continuous with an input from the previous loop or instead constituting a new input. These comparative markers are optimally built into the input-interpreting software objects. For example, a software object that initially interprets the data from the sensor array will often need to track several variables at the same time, including the reconstructed input location, the reconstructed input location from the previous loop of the software, the pressure at the present input location, the pressure at the previous input location, the numerical order in which the program register this input in relation to other simultaneous inputs, and, in the case of certain desired effects, the ‘width’ of the input—i.e. the number of sensors which contributed to the reading of the particular input location. Obviously the objects have to be coded in such a way that a set can simultaneously be interpreted. Each time the loop is run, the program can then compare each present input cluster with previous input cluster. Then once the data has been assessed at that level, the input data objects can trigger particular control effects, whether they be audio, visual, or based on movement or anything else. For example, in the example in FIG. 9 , the input data objects can be directed to send MIDI notes on and off which then produce a signal to the loud speaker 17 corresponding to a musical note. The change in the data between loops can, for example, either be used to send a new note, or to hold the existing note, increase its volume, or bend its pitch, depending on the desire sound and the comparative statements built into the code. Simply put, the program can be written such that the pitch and timbre of the musical note depends on the location on the undulating surface 12 that is pressed by the user. The volume of the note produced can be set in accordance by the magnitude of the pressure sensed by the sensor array 14 . Variations on these techniques for the purpose of other applications will be clear to a skilled programmer of control interlaces. [0074] The materials used in the keyboard are as follows: The body is made of silicone rubber, specifically Silskin 10 produced by Notcutt Ltd (Homewood Farm, Newark Lane, Ripley, Surrey GU23 6DJ) with added deadener silicone additive (also available from by Notcutt Ltd) to produce a material with a shore hardness on 00-1. The array of sensors 4 is formed from FSR strip produced by Interlink Electronics The processor 16 is an Arduino microprocessor produced by Tinker [0078] The flow diagram of FIG. 12 shows the sequence of the above-described operation of the interlace (steps 1 to 4 Å). [0079] In one embodiment, the peaks can correspond to the notes on a standard keyboard. However, it is also possible to press the trough between the peaks, which will be detected by a different combination of sensors to produce different pitches, e.g. microtonal effects, or other effects depending on the nature of the algorithms software which interprets the data. Other effects that can be produced have already been discussed above and so will not be repeated here. [0080] FIG. 9 shows a keyboard in accordance with the present invention showing not only white notes (peaks 26 ) but also black notes (peaks 25 ). In addition, touch sensitive surfaces 21 , 22 and 23 (or “sliders”) are provided above and below the keyboard that provide glissando effects when a user slides a finger, thumb or palm along it. 21 is a portamento slider, 22 is a further lower portamento slider and 23 is a palm effects slider. Indentations 24 are provided where the keys 25 , 26 meet the portamento slider 21 to provide a smooth transition surface for moving a finger smoothly (i.e. without encountering obstacles) between the main area of the keyboard (keys 25 , 26 ), and the slider 21 . There arc other ways that one can create the particular surface of the sliders depending on the precise effects and transitions one wants to achieve. [0081] FIG. 10 shows part only of the keyboard of FIG. 9 and in particular shows the peaks 27 and troughs 28 between the white and black notes. [0082] Returning now to FIG. 9 , it can be seen that the peaks can be oscillated by the finger to produce a vibrato effect or could be grasped or pinched to produce a sound that cannot be achieved using a standard keyboard. The softness of the material also means that once the surface has been depressed one can either move one's finger and slide it across the material, or one can use the flexibility of the internal layer to push the material one way or another thus giving a different reading on the underlying sensor array. Enabling this technique is important because it mean that the effects of vibrato and tremolo can be used simultaneously. [0083] Referring now to FIG. 11 , there is shown the equivalent of a computer keyboard, which has been divided into two sections, one for each hand. The top surfaces of the two keaboard parts carry elevated areas 20 corresponding to alphanumeric keys and function keys; the elevated areas can be felt by the user and so provide tactile feedback. The keyboard has a similar structure to that shown in FIG. 9 so that, when a key is pressed, a force is transmitted though soft resilient material to an array of sensors (not shown). From the combination of sensors that is triggered by a downward pressure on the top surface, it is possible to tell which key has been pressed and this detection is performed in a microprocessor, which produces an output signal causing an alphanumeric character corresponding to the pressed key to be displayed on a screen and a corresponding character to be stored into a memory. [0084] The two parts of the keyboard shown in FIG. 11 correspond in shape to the user's right and left hands. A region 34 corresponds to the place where a user will place his palm. A pointer (cursor) on the screen can be moved by the user exerting a pressure on this section 34 ; the pressure will be detected by the array of sensors and the cursor moved in accordance with the direction of the force exerted by the user in the area 34 , as detected by the signature of sensor signals. Because the material used for making the keyboard is soft, the different signatures corresponding to the different directions of the pressure exerted by the user on the region 34 can readily be detected by the sensors even if the palm of the hand does not move across the surface in the region 34 . Thus the palm can control the mouse function of directing the cursor on screen. [0085] It is also possible to provide peaks, for example peaks 32 , that can be squeezed between the user's thumb and forefinger to perform functions on the screen, for example to pick up an object displayed on the screen and drag it to a different location. [0086] It is possible to provide control, shift and command functions not only by pressing buttons but also by specific gestures performed in connection with the individual keys. [0087] These slightly more detailed examples indicate the range of applications that fall within the scope of the present invention. A person knowledgeable in the field of interface design could appreciate the breadth of possible applications that the present invention makes possible. To give just one example, although the present invention has been described in connection with user who directly touches the three-dimensional input surface, certain applications can be constructed using a similar interface which is utilized in a mechanized context, for example in a robotic context where three-dimensional force sensitive input devices can be used to create more sensitive moving joints that can sense the relative distribution of weight and can be formed to fit the mechanical parts. In these applications, the scaffolding system depicted in FIG. 6 provides an especially useful tool since it can enable movement and sensitivity in certain directions, and restrict movement in other directions, while remaining lightweight and durable. Thus a range of applications in robotics fall within the scope of the present application.

An interface is for inputting data into a processor. The interface has a three-dimensionally shaped input surface and comprises an array of sensors responsive to forces applied to the input surface and providing an input to the processor capable of registering the magnitude of the forces applied to the array of sensors and interpreting the location of pressure on the input surface, and a three-dimensionally shaped layer of soft resilient material arranged between the three-dimensionally shaped input surface and the array of sensors and capable of transmitting forces exerted on the three-dimensionally shaped input surface to the sensors.

[0001] Cross-reference is made to U.S. Utility patent application Ser. No. 12/437,576 entitled “Li-ion Battery with Selective Moderating Material” by John F. Christensen et al., which was filed on May 8, 2009; U.S. Utility patent application Ser. No. 12/437,592 entitled “Li-ion Battery with Blended Electrode” by John F. Christensen et al., which was filed on May 8, 2009; U.S. Utility patent application Ser. No. 12/437,606 entitled “Li-ion Battery with Variable Volume Reservoir” by John F. Christensen et al., which was filed on May 8, 2009; U.S. Utility patent application Ser. No. 12/437,622 entitled “Li-ion Battery with Over-charge/Over-discharge Failsafe” by John F. Christensen et al., which was filed on May 8, 2009; U.S. Utility patent application Ser. No. 12/437,643 entitled “System and Method for Pressure Determination in a Li-ion Battery” by John F. Christensen et al., which was filed on May 8, 2009; U.S. Utility patent application Ser. No. 12/437,745 entitled “Li-ion Battery with Load Leveler” by John F. Christensen et al., which was filed on May 8, 2009; U.S. Utility patent application Ser. No. 12/437,774 entitled “Li-ion Battery with Anode Coating” by Boris Kozinsky et al., which was filed on May 8, 2009; U.S. Utility patent application Ser. No. 12/437,791 entitled “Li-ion Battery with Anode Expansion Area” by Boris Kozinsky et al., which was filed on May 8, 2009; U.S. Utility patent application Ser. No. 12/437,822 entitled “Li-ion Battery with Porous Silicon Anode” by Boris Kozinsky et al., which was filed on May 8, 2009; U.S. Utility patent application Ser. No. [Attorney Docket No. 1576-0308] entitled “System and Method for Charging and Discharging a Li-ion Battery” by Nalin Chaturvedi et al., which was filed on May 8, 2009; and U.S. Utility patent application Ser. No. [Attorney Docket No. 1576-0310] entitled “System and Method for Charging and Discharging a Li-ion Battery Pack” by Nalin Chaturvedi et al., which was filed on May 8, 2009, the entirety of each of which is incorporated herein by reference. The principles of the present invention may be combined with features disclosed in those patent applications. FIELD OF THE INVENTION [0002] This invention relates to batteries and more particularly to lithium-ion batteries. BACKGROUND [0003] Batteries are a useful source of stored energy that can be incorporated into a number of systems. Rechargeable lithium-ion batteries are attractive energy storage systems for portable electronics and electric and hybrid-electric vehicles because of their high specific energy compared to other electrochemical energy storage devices. In particular, batteries with a form of lithium metal incorporated into the negative electrode afford exceptionally high specific energy (in Wh/kg) and energy density (in Wh/L) compared to batteries with conventional carbonaceous negative electrodes. [0004] When high-specific-capacity negative electrodes such as lithium are used in a battery, the maximum benefit of the capacity increase over conventional systems is realized when a high-capacity positive electrode active material is also used. Conventional lithium-intercalating oxides (e.g., LiCoO 2 , LiNi 0.8 Co 0.15 Al 0.05 O 2 , Li 1.1 Ni 0.3 Co 0.3 Mn 0.3 O 2 ) are typically limited to a theoretical capacity of ˜280 mAh/g (based on the mass of the lithiated oxide) and a practical capacity of 180 to 250 mAh/g. In comparison, the specific capacity of lithium metal is about 3863 mAh/g. The highest theoretical capacity achievable for a lithium-ion positive electrode is 1168 mAh/g (based on the mass of the lithiated material), which is shared by Li 2 S and Li 2 O 2 . Other high-capacity materials including BiF 3 (303 mAh/g, lithiated) and FeF 3 (712 mAh/g, lithiated) are identified in Amatucci, G. G. and N. Pereira, Fluoride based electrode materials for advanced energy storage devices. Journal of Fluorine Chemistry, 2007. 128(4): p. 243-262. All of the foregoing materials, however, react with lithium at a lower voltage compared to conventional oxide positive electrodes, hence limiting the theoretical specific energy. The theoretical specific energies of the foregoing materials, however, are very high (>800 Wh/kg, compared to a maximum of ˜500 Wh/kg for a cell with lithium negative and conventional oxide positive electrodes). [0005] Lithium/sulfur (Li/S) batteries are particularly attractive because of the balance between high specific energy (i.e., >350 Wh/kg has been demonstrated), rate capability, and cycle life (>50 cycles). Only lithium/air batteries have a higher theoretical specific energy. Lithium/air batteries, however, have very limited rechargeability and are still considered primary batteries. [0006] Li/S batteries also have limitations. By way of example, the United States Advanced Battery Consortium has established a goal of >1000 cycles for batteries used in powering an electric vehicle. Li/S batteries, however, exhibit relatively high capacity fade, thereby limiting the useful lifespan of Li/S batteries. [0007] One mechanism which may contribute to capacity fade of Li/S batteries is the manner in which the sulfur reacts with lithium. In general, sulfur reacts with lithium ions during battery discharge to form polysulfides (Li x S), which may be soluble in the electrolyte. These polysulfides react further with lithium (i.e., the value of x increases from ¼ to ⅓ to ½ to 1) until Li 2 S 2 is formed, which reacts rapidly to form Li 2 S. In Li/S batteries described in the literature, both Li 2 S 2 and Li 2 S are generally insoluble in the electrolyte. Hence, in a system in which intermediate polysulfides are soluble, each complete cycle consists of soluble-solid phase changes, which may impact the integrity of the composite electrode structure. [0008] Specifically, Li 2 S may deposit preferentially near the separator when the current through the depth of the positive electrode is non-uniform. Non-uniformity is particularly problematic at high discharge rates. Any such preferential deposition can block pores of the electrode, putting stress on the electronically conducting matrix and/or isolating an area from the composite electrode. All of these processes may lead to capacity fade or impedance rise in the battery. [0009] Moreover, soluble polysulfides are mobile in the electrolyte and, depending on the type of separator that is used, may diffuse to the negative electrode where the soluble polysulfides may becoming more lithiated through reactions with the lithium electrode. The lithiated polysulfide may then diffuse back through the separator to the positive electrode where some of the lithium is passed to less lithiated polysulfides. This overall shuttle process of lithium from the negative electrode to the positive electrode by polysulfides is a mechanism of self discharge which reduces the cycling efficiency of the battery and which may lead to permanent capacity loss. [0010] Some attempts to mitigate capacity fade of Li/S batteries rely upon immobilization of the sulfur in the positive electrode via a polymer encapsulation or the use of a high-molecular weight solvent system in which polysulfides do not dissolve. In these batteries, the phase change and self-discharge characteristics inherent in the above-described Li/S system are eliminated. These systems have a higher demonstrated cycle life at the expense of high rate capability and capacity utilization. [0011] In the case of a Li/S battery, however, the sulfur active material increases in volume by ˜80% as it becomes lithiated during battery discharge. Thus, an all solid-state cathode, composed of sulfur (or lithiated sulfur) and a mixed conducting material, particularly if the latter is a ceramic, is susceptible to fracture due to the volume change upon battery cycling. Fracture of the cathode can result in capacity fade and is a potential safety hazard due to venting of the cell. Other materials which exhibit desired capabilities when incorporated into a battery also exhibit significant increases in volume. By way of example, LiSi, typically used as an anode material, exhibits a large increase in volume during operation. [0012] What is needed therefore is a battery that provides the benefits of materials that exhibit large volume changes during operation of the cell while reducing the likelihood of fracture of material or internal shorts within the cell. SUMMARY [0013] In accordance with one embodiment an electrochemical cell includes a first electrode, and a second electrode spaced apart from the first electrode, the second electrode including, a current collector, an electrically conducting rigid support frame electrically connected to the current collector, and an active material coated to the rigid support frame. [0014] In accordance with another embodiment, an electrochemical cell includes a first electrode, and a second electrode spaced apart from the first electrode, the second electrode including, a current collector, an electrically conducting first support wall electrically connected to the current collector, an electrically conducting second support wall spaced apart from the first support wall and electrically connected to first support wall, and an active material coated to the first support wall and the second support wall. BRIEF DESCRIPTION OF THE DRAWINGS [0015] FIG. 1 depicts a schematic of a lithium ion cell including a cathode and an anode with a rigid framework of nanowires; [0016] FIG. 2 depicts a schematic of a lithium ion cell including a cathode and an anode with a rigid framework configured to provide directional lithium ion coating of the framework; [0017] FIG. 3 depicts a schematic of another embodiment of a lithium ion cell including a cathode and an anode with a rigid framework configured to provide directional lithium ion coating of the framework; and [0018] FIG. 4 depicts a schematic of another embodiment of a lithium ion cell including a cathode and an anode with a rigid framework configured to provide directional lithium ion coating of the framework. DESCRIPTION [0019] For the purposes of promoting an understanding of the principles of the invention, reference will now be made to the embodiments illustrated in the drawings and described in the following written specification. It is understood that no limitation to the scope of the invention is thereby intended. It is further understood that the present invention includes any alterations and modifications to the illustrated embodiments and includes further applications of the principles of the invention as would normally occur to one skilled in the art to which this invention pertains. [0020] FIG. 1 depicts a lithium-ion cell 100 , which includes a negative electrode 102 , a positive electrode 104 , and a separator region 106 between the negative electrode 102 and the positive electrode 104 . The negative electrode 102 includes a current collector 108 . A first support wall 110 is attached to the current collector 108 on one side while the other side of the support wall 1 10 faces and is spaced apart from another support wall 112 . [0021] The first support wall 110 and the second support wall 112 are formed from nanotubes, nanowires, or conducting fibers such as carbon 114 which in this embodiment are formed as a grid from a material onto which lithium 116 plates, although other materials such as graphite particles may be used. The first support wall 110 and the second support wall 112 are connected such that lithium ions can migrate between the support walls 110 and 112 . In one embodiment, a single ply of woven material is folded to provide facing surfaces of the support walls 110 and 112 . In one alternative embodiment, a solid LI-ion conductor, such as lithium phosphate, lisicon, a lithium-conducting polymer or glass, is used to create connections between support structures for lithium migration. In a further embodiment, an electrolyte is used within the electrode 102 to provide a transfer path. The support walls 110 and 112 may be spaced apart as depicted in FIG. 1 or they may be in contact with the other support wall 110 or 112 along the facing surfaces. [0022] The separator region 106 includes an electrolyte with a lithium cation and serves as a physical and electrical barrier between the negative electrode 102 and the positive electrode 104 so that the electrodes are not electronically connected within the cell 100 while allowing transfer of lithium ions between the negative electrode 102 and the positive electrode 104 . [0023] The positive electrode 104 includes a current collector 120 and an active portion 122 into which lithium can be inserted. The active portion 122 may include a form of sulfur and may be entirely sulfur. If desired, the positive electrode 104 may include a support structure similar to the support walls 110 and 112 . [0024] The lithium-ion cell 100 operates in a manner similar to the lithium-ion battery cell disclosed in U.S. patent application Ser. No. 11/477,404, filed on Jun. 28, 2006, the contents of which are herein incorporated in their entirety by reference. In general, electrons are generated at the negative electrode 102 during discharging and an equal amount of electrons are consumed at the positive electrode 104 as lithium and electrons move in the direction of the arrow 136 of FIG. 1 . [0025] In the ideal discharging of the cell 100 , the electrons are generated at the negative electrode 102 because there is extraction via oxidation of lithium ions from the lithium 116 plated on the nanowires 114 of the negative electrode 102 , and the electrons are consumed at the positive electrode 104 because metal cations or sulfur ions change oxidation state in the positive electrode 104 . During charging, the reactions are reversed, with lithium and electrons moving in the direction of the arrow 138 . [0026] As lithium ions are inserted into the active portion 122 , the volume of the active portion 122 increases. As the volume of the active portion 122 increases, the pressure within the positive electrode 104 increases. The increased pressure in the positive electrode 104 , in embodiments incorporating a fluid such as a fluid electrolyte, causes the fluid to flow toward the negative electrode 102 . Because the nanowires 114 do not fill the entire negative electrode 102 , the fluid can flow into the negative electrode 102 . The rigid support walls 110 and 112 thus provide an expansion volume within the negative electrode 102 . Additionally, the rigidity of the support structures 110 and 112 protects the active material in the electrode 102 from the volume change in the positive electrode 104 . [0027] As lithium plates onto the nanowires 114 , the lithium 116 may plate predominantly in directions toward another nanowire 114 . Specifically, since one side of the rigid wall 110 is mounted to the current collector 108 , lithium will not plate onto that surface portion. Thus, plating of lithium on the support wall 110 occurs predominantly on the sides of the nanowires 114 facing other nanowires 114 within the support wall 110 or nanowires 114 in the opposing support wall 112 . Likewise, since one side of the nanowires 114 in the rigid wall 112 is attached to the separator layer 106 , lithium will not plate onto that surface portion. Thus, plating of lithium on the support wall 112 occurs predominantly on the sides of the nanowires 114 facing other nanowires 114 within the support wall 112 or nanowires 114 in the opposing support wall 110 . [0028] Accordingly, any deformation of the lithium layer 116 on the support wall 110 will typically occur in a direction that is not directly toward the current collector 108 and any deformation of the lithium layer 116 on the support wall 110 will typically occur in a direction that is not directly toward the separator layer 106 . Thus, deformation or dendrites must extend for a significantly longer distance before any significant deleterious effects on the current collector 108 or the separator layer 106 are generated. [0029] The benefits of providing a rigid support wall that provides protection of active material from volume changes within the cell and which may promote directional plating of lithium can be increased by modifying the shape of the rigid support structure members. By way of example, FIG. 2 depicts a lithium-ion cell 200 which includes a negative electrode 202 , a positive electrode 204 , and a separator layer 206 between the negative electrode 202 and the positive electrode 204 . The negative electrode 202 includes a current collector 208 . A first support wall 210 is attached to the current collector 208 on one side while the other side of the support wall 210 faces and is spaced apart from another support wall 212 . [0030] The first support wall 210 and the second support wall 212 are formed from shaped nanowires 214 which in this embodiment are formed as a grid from a material onto which lithium 216 plates. The first support wall 210 and the second support wall 212 are connected such that lithium ions can migrate between the support walls 210 and 212 . [0031] The separator layer 206 is a lithium conductor and serves as a physical and electrical barrier between the negative electrode 202 and the positive electrode 204 so that the electrodes are not electronically connected within the cell 200 while allowing transfer of lithium ions between the negative electrode 202 and the positive electrode 204 . The positive electrode 204 includes a current collector 220 and an active portion 222 into which lithium ions can be inserted. [0032] The lithium-ion cell 200 operates in a manner similar to the lithium-ion battery cell 100 . The nanowires 214 of the lithium-ion cell 200 are formed, however, to increase the effect of directional plating as compared to the nanowires 114 of the lithium-ion cell 100 . To this end, the nanowires 214 include a mounting surface 230 attached to either the current collector 208 or to the separator layer 206 , and a plating surface 232 on which lithium is allowed to plate. The plating or active surface 232 is configured such that the plating surface 232 does not face the surface on which the nanowire 214 is mounted. Thus, plating of lithium 216 onto the nanowires 214 occurs in a direction away from the surface on which the nanowires 214 are mounted. [0033] The inclusion of a rigid framework may increase the necessary volume for a particular cell. The increased volume may be minimized by selective spacing and shaping of the members used to form the support walls. FIG. 3 , for example, depicts a lithium-ion cell 300 which includes a negative electrode 302 , a positive electrode 304 , and an electrolyte layer 306 between the negative electrode 302 and the positive electrode 304 . The negative electrode 302 includes a current collector 308 . A first support wall 310 is attached to the current collector 308 on one side while the other side of the support wall 310 faces and is spaced apart from another support wall 312 which is connected to the electrolyte layer 306 . [0034] The first support wall 310 is formed as a solid base portion 314 from which shaped protrusions 316 extend. The second support wall 312 is formed from shaped wires 318 which in this embodiment are formed into a grid which may include openings to the electrolyte layer 306 . Shaped plating surfaces 320 are supported by the wires 318 and extend toward the support wall 310 . Both the support wall 310 and the support wall 312 are formed from a material onto which lithium 322 plates. The support wall 310 and the support wall 312 may be shaped using semiconductor chip forming processes or other manufacturing processes. The first support wall 310 and the second support wall 312 are connected such that lithium ions can migrate between the support walls 310 and 312 . [0035] The electrolyte layer 306 includes an electrolyte with a lithium cation and serves as a physical and electrical barrier between the negative electrode 302 and the positive electrode 304 so that the electrodes are not electronically connected within the cell 300 while allowing transfer of lithium ions between the negative electrode 302 and the positive electrode 304 . The positive electrode 304 includes a current collector 324 and an active portion 326 into which lithium can be inserted. [0036] The lithium-ion cell 300 operates in a manner similar to the lithium-ion battery cell 200 . The support wall 310 and the support wall 312 of the lithium-ion cell 300 are formed to increase the effect of directional plating like the nanowires 214 of the lithium-ion cell 200 . The controlled manufacturing process used to form the support walls 310 and 312 , however, ensure a more uniform distance between the plating surfaces of support wall 310 and the support wall 312 and any other surface in the cell 300 . Accordingly, usefulness of the space within the electrode 302 is optimized. Moreover, the support wall 310 and the support wall 312 are formed to allow interlacing of plating surfaces to further minimize space requirements. [0037] While the foregoing embodiments incorporate wall structures which are highly ordered, a wall structure with randomly oriented support members may be used. To this end, FIG. 4 depicts a lithium-ion cell 400 , which includes a negative electrode 402 , a positive electrode 404 , and a separator region 406 between the negative electrode 402 and the positive electrode 404 . The negative electrode 402 , the positive electrode 404 , and the separator region 406 are located within a pouch 408 . [0038] The negative electrode 402 includes a support structure 410 which is in electrical contact with a current collector 412 and an electrolyte 414 . The support structure 410 is made from nanotubes or nanowires 416 which in this embodiment are a conductive mixture onto which lithium or some other active material can plate. The conductive material in the nanowires 416 may include carbon. The nanowires 416 are randomly oriented. Random orientation may be obtained, for example, by drying of a slurry incorporating a plurality of nanowires 416 . Once dried, the nanowires 416 may be compacted to reduce porosity. [0039] The separator region 406 in one embodiment includes an electrolyte with a lithium cation and serves as a physical and electrical barrier between the negative electrode 402 and the positive electrode 404 so that the electrodes are not electronically connected within the cell 400 while allowing transfer of lithium ions between the negative electrode 402 and the positive electrode 404 . [0040] The positive electrode 404 includes active material 420 into which lithium can be inserted, inert material 422 , the electrolyte 414 , and a current collector 426 . The active material 420 includes a form of sulfur and may be entirely sulfur. [0041] The lithium-ion cell 400 operates in a manner similar to the lithium-ion battery cells 100 , 200 , and 300 . Although the nanowires 416 are randomly oriented, most of the active surfaces of the nanowires 416 are either spaced away from the outer walls of the nanowire support structure 410 or oriented away from an immediately adjacent structure such as the current collector 412 . Accordingly, while providing a substantial volume occupied only by electrolyte 414 which can easily be displaced as active material plates onto the support structure 410 , any dendrite formation is less likely to establish an internal short. [0042] While the invention has been illustrated and described in detail in the drawings and foregoing description, the same should be considered as illustrative and not restrictive in character. It is understood that only the preferred embodiments have been presented and that all changes, modifications and further applications that come within the spirit of the invention are desired to be protected.

An electrochemical cell in one embodiment includes a first electrode, and a second electrode spaced apart from the first electrode, the second electrode including, a current collector, an electrically conducting rigid support frame electrically connected to the current collector, and an active material coated to the rigid support frame.

BACKGROUND OF DISCLOSURE [0001] 1. Field of the Disclosure [0002] Embodiments disclosed herein relate generally to thrusters that apply a force to a drill bit during the drilling of an underground formation. In another aspect, embodiments disclosed herein relate to control of a thrust force applied to a drill bit by a thruster. More specifically, embodiments disclosed herein relate to controlling a pressure or differential pressure across a thrust piston, thereby limiting the maximum applied thrust force. [0003] 2. Background [0004] Hydraulic thrusters are used for applying a force to an earth boring drill bit, independent of the drill string weight. Although thrusters may be used during vertical or inclined drilling, hydraulic thrusters are generally advantageous in horizontal or near-horizontal drilling. During horizontal drilling, especially in long horizontal sections, a significant portion of the weight of the drill stem is directed toward the low side of the hole, detracting from the weight available for bit thrust. Hydraulic thrusters allow for extended reach drilling, the drilling of multiple horizontal wells from a single platform, decreasing the drilling costs associated with producing reservoirs that are offshore, in arctic regions, mountains, or near large cities. [0005] The thruster is a telescoping tube arrangement that allows the drill bit to advance while the tubing string is supported in a rather stationary position at the surface. When the thruster has advanced its full stroke, or a notable portion thereof, the drill string is lowered from the surface, causing the upper end of the thruster to slide down and reset the thruster for the next stroke. When the drilling kelly or the stand being drilled down by the top drive reaches the drill rig floor, circulation is interrupted and another piece of tubing is added to the string at the surface or the coiled tubing is further unspooled into the wellbore. This drilling procedure also causes the thruster to reset. [0006] Hydraulic thrusters are described in, for example, U.S. Pat. No. 4,615,401 and patents referenced therein (U.S. Pat. Nos. 3,298,449, 3,399,738, 3,797,589, 3,799,277, 4,040,494, and 4,040,495), each of which is assigned to the assignee of the present invention, and each of which is incorporated herein by reference. In the '401 patent, the hydraulic thruster includes a mandrel and sleeve forming two expandable chambers with wall anchors annularly disposed about the sleeve responsive to a pressure differential between a chamber and the bore hole pressure. Valves and actuators are provided to extend and retract a piston between two extremes of relative axial motion between the mandrel and sleeve. [0007] Hydraulic thrusters are also described in U.S. Pat. No. 5,205,364. In the '364 patent, the hydraulic thruster includes a telescoping assembly for transmitting hydraulic force to the drill bit at the bottom of the tool. The internal hydraulic characteristics of the tool may be varied to vary the force exerted during extension and retraction of the telescoping assembly. The hydraulic characteristics are varied by varying the surface area within the drill tool on which the flow of drilling mud may act when producing the hydraulic force. [0008] Other patents describing thrusters or equipment for controlling force or weight on the bit, for example, may include U.S. Pat. Nos. 5,316,094, 6,601,652, 7,156,181, 5,476,148, 5,884,716, 5,806,611, 6003,606, 6,230,813, and 6,286,592, and U.S. Patent Application Publication No. 20010045300. [0009] Referring now to FIG. 1 , a simplified cross-sectional view of a prior art thruster 1 is illustrated. Thruster 1 , shown in the retracted position, may include an inner mandrel assembly 2 , which may include one or more tubular components. Threads 3 may connect inner mandrel assembly 2 to the lower drill stem (not shown) toward the bit (not shown). Threads 4 may connect inner mandrel assembly 2 to the upper drill stem (not shown). Inner mandrel assembly 2 is disposed in and is axially movable with respect to outer tubular assembly 5 . One or more anchor pistons 6 may be provided to anchor thruster 1 with respect to the hole wall (not shown). Drilling fluid supplied to the bore 2 A of inner mandrel 2 and to the drill bit (not shown) defines a high pressure area, and drilling fluid returning from the bit in the annulus 7 formed between the outer tubular assembly 6 and the hole wall defines a low pressure area. During thrusting, thrust mechanism 8 may be used to allow the high pressure drilling fluid into chamber A, allowing fluid in chamber B to escape to annulus 7 , thus creating a pressure differential across thrust mechanism 8 , causing the inner mandrel 2 to advance in direction 9 , and putting weight on the bit corresponding to the thrust force generated by the pressure differential. [0010] A cross-sectional view of a simplified thrust mechanism 8 , which may be used in the thruster of FIG. 1 , is illustrated in FIG. 2 . Thrust mechanism 8 may include an inner tubular member 12 and an outer tubular member 14 . Drilling mud flowing through the bore 16 of inner tubular member 12 flows to the drill bit (not shown), and returns to the surface via annulus 18 , such as between outer tubular member 14 and a drill casing (not shown). When mud is flowing through thruster 1 ( FIG. 1 ), bore 16 is at a higher pressure than fluid returning through annulus 18 . A piston 20 , separating a first fluid chamber 22 and a second fluid chamber 24 , may transmit an axial force 26 to inner tubular member 12 . During thrusting, high pressure drilling mud flows from the bore 16 of the thruster 1 through inlet 28 into first fluid chamber 22 , displacing fluid in second fluid chamber 24 through outlet 30 and causing the inner tubular member 12 to advance in the direction of axial force 26 . The axial force 26 that is generated, for example, may be a function of the differential pressure between the fluid in bore 16 and annulus 18 . [0011] Many of the patents cited above use such a differential pressure to control the force applied to the drill bit. Although not shown in FIG. 2 , thrust mechanism 8 may typically include ball valves, springs, and other mechanisms to control the flow of fluid into and from the high and low pressure chambers, respectively, during thrusting and retraction. One problem associated with this type of thruster technology includes the need to estimate the pressure and required thrust force prior to drilling. The thruster and the associated internal parts are generally selected and fabricated at the surface prior to installation on a drill string, and many of the parts used to control fluid flow, such as springs, check valves, flow orifices, and others, are sized and selected based on an expected downhole pressure. [0012] Often, an actual downhole pressure differs from the expected downhole pressure. The difference between actual and expected downhole pressure, even by as little as 25 or 50 psi, may result in ineffective control of the force applied to the drill bit by the thruster, often as a result of the thrust mechanism fully opening or fully closing. Additionally, fluctuations in pressure drop across the bit and changes in the weight of the drilling fluid used (and hence bore pressure) may also result in ineffective control of the force applied to the drill bit by the thruster. Ineffective thruster control may lead to stalls, motor wear, stuck bits, and inefficient rate of penetration, among other problems known to those skilled in the art. [0013] Various methods and apparatus have been proposed to compensate for a change in downhole conditions and to minimize the effect such changes have on the operation of the thruster and the force applied to the drill bit. For example, a pressure-modulation valve assembly is disclosed in U.S. Pat. No. 6,102,138. Such methods and apparatus unnecessarily increase the total number of drilling components of a drill string, where the additional apparatus may be prone to failure or malfunction due to various conditions encountered during drilling. [0014] Accordingly, there exists a need for a thruster that may control the force applied to a drill bit independent of the downhole pressure or the pressure drop across the motor and bit. Additionally, there exists a need for a thruster that may control the force applied to the bit independent of the pressure of the supplied drilling fluid. SUMMARY OF THE DISCLOSURE [0015] In one aspect, embodiments disclosed herein relate to a drilling system, including: a drill bit; and a thruster to apply a force to the drill bit. The thruster may include: an inner tubular member disposed within and configured to axially move within an outer tubular member; a thrust piston to transmit a hydraulic force to the inner tubular member, the thrust piston separating an upstream fluid chamber and a downstream fluid chamber between the inner and outer tubular members; at least one pressure switch fluidly connected to the downstream fluid chamber to control flow of a fluid to and from the downstream fluid chamber via at least one fluid inlet and at least one fluid outlet. [0016] In another aspect, embodiments disclosed herein relate to a thruster, including: an inner tubular member disposed within and configured to axially move within an outer tubular member; a thrust piston to transmit a hydraulic force to the inner tubular member, the thrust piston separating an upstream fluid chamber and a downstream fluid chamber between the inner and outer tubular members; a downstream valve member mechanically coupled to a downstream magneto-actuator and disposed in the downstream fluid chamber; and at least one pressure switch fluidly coupled to the downstream fluid chamber to control a position of the downstream valve member via the magneto-actuator; wherein the position of the downstream valve member affects a flow of a fluid to and from the downstream fluid chamber via at least one fluid inlet and at least one fluid outlet. [0017] In another aspect, embodiments disclosed herein relate to a process to drill an underground formation. The process may include: supplying a fluid to a thruster, wherein the thruster includes: an inner tubular member disposed within and configured to axially move within an outer tubular member; a thrust piston to transmit a hydraulic force to the inner tubular member, the piston separating an upstream fluid chamber and a downstream fluid chamber between the inner tubular member and the outer tubular member; at least one pressure switch fluidly connected to the downstream fluid chamber; and regulating a flow of the fluid to and from the downstream fluid chamber in response to a signal from the at least one downstream pressure switch to maintain the hydraulic force applied to the inner tubular member proximate a hydraulic force set point. [0018] Other aspects and advantages will be apparent from the following description and the appended claims. BRIEF DESCRIPTION OF DRAWINGS [0019] FIG. 1 is a simplified schematic drawing of a prior art thruster. [0020] FIG. 2 is a schematic drawing of a simplified thrust mechanism that may be used in the prior art thruster of FIG. 1 . [0021] FIG. 3 is a simplified schematic drawing of a thruster according to embodiments disclosed herein. [0022] FIG. 3A is a simplified schematic drawing of an actuator useful in embodiments of the thrusters described herein. [0023] FIG. 3B is a simplified schematic drawing of an actuator useful in embodiments of the thrusters described herein. [0024] FIG. 3C is a simplified schematic drawing of an actuator useful in embodiments of the thrusters described herein. [0025] FIG. 4 is a simplified schematic drawing of a thruster according to embodiments disclosed herein. [0026] FIG. 5 is a simplified schematic drawing of a thruster according to embodiments disclosed herein. [0027] FIG. 6 is a simplified schematic drawing of a thruster according to embodiments disclosed herein. DETAILED DESCRIPTION [0028] In one aspect, embodiments disclosed herein relate to control of a thrust force applied to a drill bit by a thruster. More specifically, embodiments disclosed herein relate to controlling a pressure or differential pressure across a thrust piston, thereby limiting the maximum applied thrust force. Other embodiments disclosed herein relate to a method of drilling a formation using a thruster that may limit the thrust force applied to the bit independent of bore and annulus fluid pressures. [0029] As described above, prior art thrusters generate an axial force based upon a difference in bore and annulus pressures. In contrast, thrusters disclosed herein include mechanisms to regulate the pressure in one or both of the upstream and downstream fluid chambers. The axial force generated according to embodiments disclosed herein, for example, may be a function of the differential pressure between the fluid in the upstream and downstream fluid chambers. [0030] Referring now to FIG. 3 , a simplified schematic drawing of a thruster 50 according to embodiments disclosed herein is illustrated. Thruster 50 may include an inner tubular member 52 and an outer tubular member 54 . Drilling mud flowing through the bore 56 of inner tubular member 52 flows to the drill bit (not shown), and returns to the surface via annulus 58 , such as between outer tubular member 54 and a drill casing (not shown). When mud is flowing through thruster 50 , bore 56 is at a higher pressure than fluid returning through annulus 58 . A thrust piston 60 , separating an upstream fluid chamber 62 and a downstream fluid chamber 64 , may transmit an axial force 66 to inner tubular member 52 . During thrusting, high pressure drilling mud flows from the bore 56 of the thruster 50 through inlet 68 into upstream fluid chamber 62 , displacing low pressure fluid in downstream fluid chamber 64 through outlet 70 and causing the inner tubular member 52 to advance in the direction of axial force 66 . [0031] To regulate thrust force, or differential pressure between the upstream chamber 62 and the downstream chamber 64 , for example, thruster 50 may include a pressure switch 72 , which may be in fluid communication with the downstream fluid chamber 64 . Pressure switch 72 , in some embodiments, may be a pressure limit switch, activating at a pressure set point. When the fluid in chamber 64 reaches a predetermined set point pressure, the pressure switch 72 may actuate. Upon actuation, pressure switch 72 may send an electronic signal to a control mechanism (not shown) for regulating the flow of fluid into or out of downstream fluid chamber 64 through downstream inlet 74 and outlet 70 . [0032] By sending a signal to regulate the flow of fluid into and out of downstream fluid chamber 64 , pressure switch 72 may limit the thrust force applied to the drill bit, thus avoiding the full on or full off scenarios often encountered with prior art thrusters. For example, by limiting the flow of fluid through outlet 70 , pressure will build in downstream fluid chamber 64 , limiting the applied thrust force. As another example, by allowing fluid to flow in through inlet 74 , pressure will also increase in downstream fluid chamber 64 , due to high pressure fluid in bore 56 , limiting the applied thrust force. [0033] The control mechanism may in turn send a signal or a current to a valve member 76 to regulate the flow of fluid into and out of downstream fluid chamber 64 . Valve member 76 may include, for example, an actuator 78 , a drive rod 80 , and a gate member 82 . The signal or current transmitted to valve member 76 may cause actuator 78 to extend or contract, as illustrated by the arrows, causing a similar displacement in drive rod 80 , causing gate 82 to open and/or close fluid inlet 74 and/or fluid outlet 70 . Other means of regulating fluid flow using a signal from a pressure switch are also contemplated herein. [0034] Actuator 78 may include any one of several types of actuators responsive to electronic signals or currents. For example, actuator 78 may include magnetostrictive actuators, shape memory alloy actuators, and linear motor actuators. Examples of each of these are illustrated in FIGS. 3A-3C . [0035] As illustrated in FIG. 3A , actuator 78 may include a magnetostrictive actuator, including permanent magnets 84 , drive rod 85 , coil 86 , preload springs 87 , and output rod 88 . Upon application of a current through coil 86 , drive rod 85 may expand or contract in response to the magnetic field generated, thereby displacing output rod 88 to control the position of the gate member 82 and thus control the flow of fluid to and from the downstream cavity 64 . [0036] As illustrated in FIG. 3B , actuator 78 may include a shape memory alloy actuator, including shape memory alloy spring 90 , piston 92 , and drive rod 94 . Upon application of an electrical current, shape memory alloy spring may expand or contract, thereby displacing piston 92 and drive rod 94 to control the position of the gate member 82 , and thus control the flow of fluid to and from the downstream cavity 64 . [0037] As illustrated in FIG. 3C , actuator 78 may include a linear motor actuator, including a stationary member 96 , a motive member 97 , and a drive rod 98 . Linear motor actuators may include flat linear motor actuators and, as illustrated, tubular linear motor actuators. In some embodiments, a signal sent from the control mechanism to the linear motor actuator may control the position of the motive member 97 , and thus drive rod 98 , with respect to stationary member 96 . In other embodiments, a signal sent from the control mechanism to the linear motor actuator may control an output force exerted on drive rod 98 . In this manner, the linear motor actuator may control the position of gate member 82 , and thus control the flow of fluid to and from the downstream cavity 64 . [0038] Referring now to FIG. 4 , a simplified schematic drawing of a thruster 100 according to embodiments disclosed herein is illustrated. Thruster 100 may include an inner tubular member 102 and an outer tubular member 104 . Drilling mud flowing through the bore 106 of inner tubular member 102 flows to the drill bit (not shown), and returns to the surface via annulus 108 , such as between outer tubular member 104 and a drill casing (not shown). When mud is flowing through thruster 100 , bore 106 is at a higher pressure than fluid returning through annulus 108 . A thrust piston 110 , separating an upstream fluid chamber 112 and a downstream fluid chamber 114 , may transmit an axial force 116 to inner tubular member 102 . During thrusting, high pressure drilling mud flows from the bore 106 of the thruster 100 through inlet 118 into upstream fluid chamber 112 , displacing low pressure fluid in downstream fluid chamber 114 through outlet 120 and causing the inner tubular member 102 to advance in the direction of axial force 116 . [0039] To regulate thrust force, or differential pressure between the upstream chamber 112 and the downstream chamber 114 , for example, thruster 100 may include a pressure switch 122 , which may be in fluid communication with each of the upstream fluid chamber 112 and the downstream fluid chamber 114 . Pressure switch 122 , in some embodiments, may be a differential pressure limit switch, activating at a differential pressure set point. When the differential pressure of the fluid in upstream and downstream chambers 112 , 114 reaches a pre-determined differential pressure set point, the pressure switch 122 may actuate. Upon actuation, pressure switch 122 may send an electronic signal to a control mechanism (not shown) for regulating the flow of fluid into or out of downstream fluid chamber 114 through downstream inlet 124 and outlet 120 . [0040] By sending a signal to regulate the flow of fluid into and out of downstream fluid chamber 114 , pressure switch 122 may regulate the thrust force applied to the drill bit, as described above. [0041] The control mechanism may in turn send a signal or a current to a valve member 126 to regulate the flow of fluid into and out of downstream fluid chamber 114 . Valve member 126 may include, for example, an actuator 128 , a drive rod 130 , and a gate member 132 . The signal or current transmitted to valve member 126 may cause actuator 128 to extend or contract, as illustrated by the arrows, causing a similar displacement in drive rod 130 , causing gate 132 to open and/or close fluid inlet 124 and/or fluid outlet 120 . [0042] Referring now to FIG. 5 , a simplified schematic drawing of a thruster 150 according to embodiments disclosed herein is illustrated. Thruster 150 may include an inner tubular member 152 and an outer tubular member 154 . Drilling mud flowing through the bore 156 of inner tubular member 152 flows to the drill bit (not shown), and returns to the surface via annulus 158 , such as between outer tubular member 154 and a drill casing (not shown). When mud is flowing through thruster 150 , bore 156 is at a higher pressure than fluid returning through annulus 158 . A thrust piston 160 , separating an upstream fluid chamber 162 and a downstream fluid chamber 164 , may transmit an axial force 166 to inner tubular member 152 . During thrusting, high pressure drilling mud flows from the bore 156 of the thruster 150 through inlet 168 into upstream fluid chamber 162 , displacing low pressure fluid in downstream fluid chamber 164 through outlet 170 and causing the inner tubular member 152 to advance in the direction of axial force 166 . [0043] To regulate thrust force, or differential pressure between the upstream chamber 162 and the downstream chamber 164 , for example, thruster 150 may include a pressure switch 172 , which may be in fluid communication with the downstream fluid chamber 164 . Pressure switch 172 , in some embodiments, may be a pressure limit switch, activating at a pressure set point. When the fluid in chamber 164 reaches a pre-determined set point pressure, the pressure switch 172 may actuate. Upon actuation, pressure switch 172 may send an electronic signal to a control mechanism (not shown) for regulating the flow of fluid into or out of downstream fluid chamber 164 through downstream inlet 174 and outlet 170 . Thruster 150 may also include a pressure switch 173 , which may be in fluid communication with the upstream fluid chamber 162 . When the fluid in chamber 162 reaches a pre-determined set point pressure, the pressure switch 173 may actuate, sending an electronic signal to a control mechanism (not shown) for regulating the flow of fluid into or out of upstream fluid chamber 162 through upstream inlet 168 and upstream outlet 175 . By sending a signal to regulate the flow of fluid into and out of upstream fluid chamber 162 and downstream fluid chamber 164 , pressure switches 173 , 172 may each, separately or collectively, limit the thrust force applied to the drill bit. [0044] The control mechanism may in turn send a signal(s) or a current(s) to valve members 176 , 177 to regulate the flow of fluid into and out of one or both of upstream and downstream fluid chambers 162 , 164 . Valve members 176 , 177 may include, respectively, for example, actuators 178 , 179 , drive rods 180 , 181 , and gate members 182 , 183 . The signal(s) or current(s) transmitted to valve members 176 , 177 may cause actuators 178 , 179 to extend or contract, as illustrated by the arrows, causing a similar displacement in drive rods 180 , 181 , causing gates 182 , 183 to open and/or close fluid inlets 174 , 175 and/or fluid outlets 170 , 171 . In some embodiments, valve action on both sides of the thrust piston 160 is required in order to have hydraulic volume flow in the upstream and downstream chambers 162 , 164 . [0045] Referring now to FIG. 6 , a simplified schematic drawing of a thruster 200 according to embodiments disclosed herein is illustrated. Thruster 200 may include an inner tubular member 202 and an outer tubular member 204 . Drilling mud flowing through the bore 206 of inner tubular member 202 flows to the drill bit (not shown), and returns to the surface via annulus 208 , such as between outer tubular member 204 and a drill casing (not shown). When mud is flowing through thruster 200 , bore 206 is at a higher pressure than fluid returning through annulus 208 . A thrust piston 210 , separating an upstream fluid chamber 212 and a downstream fluid chamber 214 , may transmit an axial force 216 to inner tubular member 202 . During thrusting, high pressure drilling mud flows from the bore 206 of the thruster 200 through inlet 218 into upstream fluid chamber 212 , displacing low pressure fluid in downstream fluid chamber 214 through outlet 220 and causing the inner tubular member 202 to advance in the direction of axial force 216 . [0046] To regulate thrust force, or differential pressure between the upstream chamber 212 and the downstream chamber 214 , for example, thruster 200 may include a differential pressure switch 222 , which may be in fluid communication with each of the upstream fluid chamber 212 and the downstream fluid chamber 214 . When the differential pressure of the fluid in upstream and downstream chambers 212 , 214 reaches a pre-determined differential pressure set point, the pressure switch 222 may actuate, sending an electronic signal to a control mechanism (not shown) for regulating the flow of fluid into or out of one or both of upstream and downstream fluid chambers 212 , 214 , thereby limiting the thrust force applied to the drill bit. [0047] The control mechanism may in turn send a signal(s) or a current(s) to valve members 226 , 227 to regulate the flow of fluid into and out of one or both of upstream and downstream fluid chambers 212 , 214 . Valve members 226 , 227 may include, respectively, for example, actuators 228 , 229 , drive rods 230 , 231 , and gate members 222 , 223 . The signal(s) or current(s) transmitted to valve members 226 , 227 may cause actuators 228 , 229 to extend or contract, as illustrated by the arrows, causing a similar displacement in drive rods 230 , 231 , causing gates 232 , 233 to open and/or close fluid inlets 224 , 225 and/or fluid outlets 220 , 221 . In some embodiments, valve action on both sides of the thrust piston 210 is required in order to have hydraulic volume flow in the upstream and downstream chambers 212 , 214 . [0048] As described above, operation and control of the thrusters described herein may be affected by remote signals, such as by actuating valves and other thruster components. In some embodiments, the control settings for the valves, actuators, and pressure switches may be adjusted using remote signals. [0049] In other embodiments, the operation and control of the thrusters described herein may be affected by down-linking a signal from the surface. For example, a signal from the surface may be used to communicate with the thruster control mechanism, such as to influence the forward movement of the thruster to initiate a change in drilling rate, a change in drilling direction, or other drilling parameters, for example. Down-linking signals, in some embodiments, may include a change in pump pressure at the surface held for a given period of time. In other embodiments, down-linking signals may include a positive and/or negative pressure pulses, such as may be actuated by a change in standpipe pressure, for example. In this manner, down-linking may be used to accurately position a well and improve drilling performance. [0050] Embodiments disclosed herein may include one or more pressure switches and/or differential pressure switches to result in the desired thrust control. In some embodiments, the pressure switches may actuate upon increasing pressure or pressure differential. In other embodiments, the pressure switches may actuate upon decreasing pressure or pressure differential. In yet other embodiments, combinations of pressure switches actuating upon increasing and decreasing pressure differential may be used, such as where a valve member opens upon increasing pressure differential in response to a signal from a first pressure switch, and the valve member closes upon decreasing pressure differential in response to a signal from a second pressure switch. Additionally, embodiments may include pressure switches and differential pressure switches in fluid communication with one or more of the upstream chamber, the downstream chamber, the inner tubular member bore, and the annulus between the outer tubular member and the hole wall, with the pressure switch actuating upon a give pressure or pressure differential so as to regulate thrust force. [0051] As described above, use of pressure switches and actuators may provide for passive thrust force control. For example, a pressure switch may actuate at a minimum or maximum desired thrust force, sensing a fully opened or fully closed condition, and thereafter adjusting the pressures in the upstream and downstream chambers. [0052] Embodiments disclosed herein may include one or more actuators to result in active thrust force control. In some embodiments, two or more actuators, of the same or different type, may be used in parallel, such as operating two or more gate members, or in series, such as to achieve a longer stroke length. Additionally, intermediate components may be used intermediate drive rod and gate member, such as lever arms and bell cranks, among others, so as to result in the desired valve action or stroke length. [0053] Embodiments disclosed herein may include two or more actuators and pressure switches in parallel to control fluid flow into and from a fluid chamber. In some embodiments, the two or more pressure switches may include different pressure set points, such that a valve member may reset prior to a subsequent cycle, for example. Pressure set points may be varied minimally so as to maintain a similar maximum thrust force upon actuation of the various switch/actuator/valve combinations. [0054] As described above, use of pressure switches and actuators in parallel or series may provide for active thrust force control. For example, when approaching a fully opened or fully closed condition, the pressure switches may actuate, adjusting the pressures in the upstream and downstream chambers and thereby operating the thruster within a desired range of thrust force. [0055] In other embodiments, two or more actuators and pressure switches may be used in series to control fluid flow into and from a fluid chamber. For example, two or more pressure switches may include different set points, such that actuators extend or contract at different pressure set points. Upon actuation of a first pressure switch/actuator pair, a minimal flow opening may be provided to limit thrust force. If differential pressure continues to increase following actuation of the first pressure switch/actuator pair, a second and subsequent pressure switch/actuator pairs may provide additional flow area to limit the thrust force applied to the bit. In this manner, thrust force may vary less significantly than an on/off type actuator/valve member. [0056] Although described with reference to the pressure chambers, one skilled in the art will recognize that embodiments of thrusters disclosed herein may include components that may be typically included in thrusters, such as the thrusters described in U.S. Pat. No. 4,615,401 and others mentioned above. For example, thrusters disclosed herein may include anchor assemblies, ball valves, seals, springs and spring assemblies, threaded connections, spacers, snap rings, bearings, pins, valve seats, and rods, among others. Components used to regulate fluid flow during resetting of the thruster may also be included. [0057] In some embodiments, additional measurement and control devices may also be used to limit or control the thrust force. For example, a sensor measuring rate of penetration may be used to actuate the valve members, thereby controlling the flow of fluid into and from the upstream and downstream fluid chambers. In this manner, rate of penetration may be maintained within a desired range, such as within an optimal range for a particular drill bit. Stroke measurement devices or position sensors may also be used to indicate the thruster position, thereby allowing an operator to slow the rate of thrust toward the end of a stroke. [0058] In some embodiments, power and currents supplied to the control mechanisms, pressure switches, and actuators may include electrical currents supplied from batteries. In other embodiments, power and currents may be supplied to the control mechanisms, pressure switches, and actuators may include electrical currents supplied from downhole generators, such as turbine generators and the like. [0059] Advantageously, embodiments disclosed herein may provide for improved thrust force control, or improved control of the weight on bit. Actuators, pressure switches and valve members described herein may advantageously limit the pressure differential between upstream and downstream chambers, thus limiting the thrust force transmitted by the thrust piston to the inner tubular member. Additionally, for embodiments limiting the pressure or pressure differential within each fluid chamber, the maximum thrust force applied may be controlled independent of fluid pressure in the inner bore and the annulus. Embodiments disclosed herein, through limiting applied thrust force, may advantageously maintain weight on bit within a desired range, improving rates of penetration, and decreasing motor wear and the occurrence of stuck bits and stalls, among other common problems known in the art. [0060] While the disclosure includes a limited number of embodiments, those skilled in the art, having benefit of this disclosure, will appreciate that other embodiments may be devised which do not depart from the scope of the present disclosure. Accordingly, the scope should be limited only by the attached claims.

A drilling system, including: a drill bit; and a thruster to apply a force to the drill bit. The thruster may include: an inner tubular member disposed within and configured to axially move within an outer tubular member; a thrust piston to transmit a hydraulic force to the inner tubular member, the thrust piston separating an upstream fluid chamber and a downstream fluid chamber between the inner and outer tubular members; at least one pressure switch fluidly connected to the downstream fluid chamber to control flow of a fluid to and from the downstream fluid chamber via at least one fluid inlet and at least one fluid outlet.

PRIOR APPLICATION INFORMATION [0001] This application claims the benefit of U.S. Provisional Application 60/579,247, filed Jun. 15, 2004 and U.S. Provisional Application 60/698,862, filed Sep. 13, 2004. BACKGROUND OF THE INVENTION [0002] A new coronavirus that caused severe acute respiratory syndrome (SARS) was identified in early 2003, and subsequently named SARS coronavirus (SARS-CoV). The virus has a high tendency to spread among humans, and the mortality can be as high as 10-15% [1,2]. The complete understanding of pathogenesis of SARS remains tentative: a recent histological study using SARS-CoV infected patient lung samples found that diffuse alveolar damage may play an important role in the progression of the disease [3]. Even though there was a significant morbidity drop this year, the likelihood of the evolution of SARS-CoV in humans and animals may result in a re-emergence of the deadly virus. [0003] Coronaviruses are enveloped viruses with single-stranded positive-sense RNA genomes typically of approximately 30 kb [4]. Most viruses in the coronavirus family cause diseases in animals, while a few, such as HCoV-229E, HCoV 0043, HCoV-NL-63 and SARS-CoV, are human pathogens [5,6,7]. Among all human coronaviruses, SARS-CoV is the only one that causes severe clinical consequences. Sequence comparisons of SARS-CoV genome sequences from different patient isolates revealed high homology; however, the sequence differences between SARS-CoV and other coronaviruses are significant. The SARS-CoV genome has about 15 predicted open reading frames, six of which can be linked to other known coronavirus genes. These genes are: 1a-1b, spike (S), envelope (E), membrane (M) and nucleocapsid (N), which was found to be a multimeric form and to be involved in host cell signal transduction regulations [8,9]. The rest of the ORFs may encode genes with as yet unknown functions [1]. [0004] The development of antivirals for SARS-CoV has been vigorously pursued after the identification of the virus with mixed successes and challenges. Tan et al screened 19 clinically approved compounds for anti SARS-CoV activity, including nucleoside analogs, interferons, protease inhibitors, reverse transcriptase inhibitors and neuraminidase inhibitors. IFNs β-1b, α-n1, α-n3 and ribavarin showed anti-viral activities at high concentrations. However, significant cytotoxic effects or lack of efficacy were also observed [10]. For instance, in two other independent assays, ribavarin was shown to have little effect against SARS-CoV replication [11,12]. IFNs α, β-1b, α-n1 and α-n3 have also been tested for their anti SARS-CoV activities [10,12], however, the moderate inhibition effect of SARS-CoV replication by interferons could only be observed at very high concentrations [10,12]. As for new drug development, glycyrrhizin was reported to possess anti-SARS-CoV activity at high concentrations [12]. Clearly, continuing to search for potent anti-SARS-CoV compounds is absolutely necessary. [0005] Unlike many other RNA viruses, coronaviruses synthesize multiple subgenomic mRNA fragments, with each subgenomic RNA usually encoding only one protein [4]. As a consequence, the transcription of coronavirus RNA is very important for virus replication. In this study, we studied the antiviral effects of ATA against SARS-CoV replication in Vero cells and found that ATA drastically inhibited virus replication by as much as 1000-fold compared to untreated controls, with little toxicity observed to be associated with ATA treatment. Anti-viral selectivity of ATA was demonstrated by its failure to inhibit adenovirus replication. Importantly, we found that ATA is a much more potent anti-SARS-CoV compound than IFNs α and β. [0006] Vaccinia virus belongs to the Poxviridae family of double stranded DNA viruses, with a genome size of approximately 190 Kbp. Poxviruses are characterized by their complex structure and large genomes, permitting them a relatively high degree of independence from their host cells and the ability to synthesize close to 200 proteins. These viruses replicate in the cytoplasm of infected cells, in discrete locations termed viral factories, which have been shown to be free of host cell organelles (50). Since the virus does not enter the nucleus, it must bode for all proteins necessary for viral transcription and replication. Vaccinia virus encodes two protein kinases, F10L and B1R, as well as a protein phosphate, H1L. These three proteins perform essential functions in the virus lifecycle, suggesting that regulation of phosphorylation/dephosphorylation events play a key role during infection. Phosphorylation acts as a mediator of signaling pathways, activating or deactivating proteins involved in transducing signals. Viruses have evolved the ability to alter host signaling to create a favorable environment for their replication. [0007] The vaccinia virus open reading frame designated H1L encodes a 19.3 Kdal dual specificity protein phosphatase that is expressed during the late stage of infection. This enzyme is transported into host cells at approximately 200 molecules per virus particle and a fraction of it is released into the cytoplasm upon uncoating (38, 47). The F10 kinase is also a virion component. This suggests that these proteins may have an immediate role in infection by regulating the phosphorylation state of specific proteins. The gene encoding H1L is well conserved amongst the Poxviruses, with homologues identified in variola virus, ectromelia virus, monkeypox virus, cowpoxvirus, myxoma virus and shope fibroma virus (38, 51). This enzyme shares the active site motif, HCXXXXXRS, common to the dual specific and protein tyrosine phosphatases (33, 34, 38). Two human homologues, termed VHR and VHX, also exist. H1L was the first phosphatase discovered with the ability to dephosphorylate both Tyr and Ser/Thr residues (37). A conserved Cysteine at position 110 within H1L has been shown to be essential for enzymatic activity towards both Tyr and Ser/Thr phosphorylated proteins. This suggests that the dephosphorylation reaction likely proceeds by a conserved mechanism in both cases (34). [0008] H1L is essential for virus viability, as evidenced by the absence of viral gene transcription when the expression of H1L is repressed (47). Since the cascade-like nature of Poxvirus transcription relies upon early transcription to direct intermediate and ultimately late gene transcription (27), this block will prove fatal to the virus. Repression of H1L also leads to hyper-phosphorlyation of several viral proteins, including the products of the F18, A14 and A17 genes (35, 49). The function of the phosphate groups on these protein remains unknown. The exact role of H1L during infection has yet to be elucidated, but besides mediating viral transcription, the phosphatase has also been shown to alter host cell signaling pathways. Phosphorylated Stat1 can be dephosphorylated by H1L in vivo, blocking the expression of IFN-y induced genes (52). This may represent one means by which vv overcomes host defenses. In this study we demonstrate that ATA can inhibit the activity of H1L in vitro, and down-regulate the Erk signaling cascade. These events are proposed to be at least two of the ways that ATA exerts its anti-viral effect. [0009] As will be seen, the aromatic polyanion Aurintricarboxylic Acid (ATA) has been shown to have a number of diverse activities although the mechanism by which ATA exhibits these effects is often poorly understood. It was initially postulated that ATA would inhibit the association of any nucleic acid binding protein with nucleic acid. Subsequent research has shown that nucleases appear to be more sensitive to inhibition by ATA than other enzymes. [0010] Specifically, ATA has been shown to inhibit the RNA transcription of vesicular stomatitis virus [13]. It has alio been shown that ATA could interact with ribosomal proteins in vitro and inhibit protein synthesis [14,15]. ATA is also believed to promote cell survival and proliferation by activation of the IGF-IR signalling pathway (Beery et al., 2001, Endocrinology 142: 3098-3107) although it was previously believed that this activity was a result of ATA's inhibition of cellular endonucleases, discussed above. [0011] ATA is also thought to inhibit transcription of iNOS genes (Tsi et al., 2002, Mol Pharmacol 101: 90-101) possibly by inhibiting upstream signal kinases. [0012] Andrew et al. (1999, Immunopharmacology 41: 1-10) suggested that ATA at a concentration of 25 μM had effects on protein phosphorylation, in addition to inhibiting endonuclease activity. [0013] Nakane et al. (1988, Eur. J. Biochem. 177: 91-96) showed that both Evans blue and ATA exhibited inhibitory effects on the in vitro activity of all DNA polymerases, including human DNA polymerases α, β, γ, DNA primase, calf-thymus terminal deoxynucleotidyltransferase, RLV reverse transcriptase, E. Coli DNA polymerase I and RNA polymerase; however, these compounds did not inhibit DNA, RNA or protein synthesis in intact cells at the concentrations which proved inhibitory in vitro, suggesting that the polymerases existed in an organized state in the nucleus, which protected them from these compounds. Similarly, Thompson and Reed (1005, Toxicol Lett 81: 141-149) showed that ATA inhibited a wide range of NAD(H)/NADP(H)-requiring enzymes in in vitro incubations using purified enzymes but the inhibitory effects were markedly reduced in incubations which more closely resembled a cellular milieu. [0014] Cushman and Sherman (1992, Biochem Biophys Res Commun 185: 85-90) showed that ATA acted as an inhibitor of HIV-1 integration protein (IN). Similarly, Schols et al (1989, PNAS 86: 3322-3326) believed that ATA was targeting the CD4 receptor and thereby interfering with HIV infection. They further noted that ATA had no inhibitory effect at subtoxic concentrations for viruses that did not require the CD4 receptor to infect cells (herpes simplex virus, cytomegalovirus and vesicular stomatitis virus specifically were tested) indicating that “ATA is not a selective inhibitor of viruses other than HIV”. [0015] Givens and Manly (NAR 3:405-418) tested the effect of ATA on RNA dependent DNA polymerases. It was believed that ATA was a nonspecific inhibitor of nucleotide-requiring enzymes in vitro but this same effect was not found in vivo. One would conclude that ATA would have a nonspecific inhibitory effect on all polymerases in a purified system, but no discernable effect under in vivo-like conditions. Furthermore, ATA was not tested against RNA dependent RNA polymerases, which is only found in RNA viruses. [0016] It is of note that the literature lists several compounds which have been reported to be equivalent to ATA under certain conditions, including Evans Blue, suramin, and polyethylene sulfonic acid. In addition, Liang et al (JBC 278: 41734-41741) also lists a number of other compounds. SUMMARY OF THE INVENTION [0017] According to a first aspect of the invention, there is provided a method of treating an individual infected with or suspected of being infected with an ATA-sensitive virus comprising administering to an individual in need of such treatment an effective amount of aurintricarboxylic acid (ATA) or a derivative thereof. [0018] According to a second aspect of the invention, there is provided a method of identifying an organism inhibited by ATA comprising: [0019] searching a protein structure database for a peptide of interest having a region homologous to R binding region of SARS-CoV RdRp; and [0020] determining if the homologous region contains catalytic residues for the protein of interest. [0021] According to a third aspect of the invention, there is provided a method of screening organisms of interest for inhibition by ATA comprising: [0022] incubating an organism of interest under appropriate growth conditions in the presence of ATA; and [0023] determining if growth of the organism of interest has been inhibited. [0024] According to a fourth aspect of the invention, there is provided a method of inhibiting growth of an organism comprising administering an effective amount of ATA or a derivative thereof, wherein the organism comprises an essential protein having a region homologous to R binding region of SARS CoV RdRp, wherein said homologous region comprises at least one catalytic residue of the protein. [0025] According to a fifth aspect of the invention, there is provided a method of preparing a medicament for treating an individual infected with or suspected of being infected with an ATA-sensitive virus comprising combining an effective amount of aurintricarboxylic acid (ATA) and a suitable excipient. BRIEF DESCRIPTION OF THE DRAWINGS [0026] FIG. 1 . Vero cells were infected with SARS-CoV and treated with dilutions of aurintricarboylic acid. At 24 h. post infection, supernatant samples were harvested for a plaque assay. The virus titres of ATA treated and untreated samples were calculation and represented by plaque-forming unit (PFU/ml). The experiments were repeated at least 3 times, with SD being approximately 10%. [0027] FIG. 2 . Vero cells were infected with SARS-CoV and treated with serially diluted concentrations of aurintrucarboxylic acid. After 24 h., cells were harvested and subjected to 4-12% SDS-PAGE; protein samples were subsequently transferred to PVDF membrane and probed with a mouse monoclonal antibody against SARS-CoV spike protein and anti-α-actinin antibodies. A rabbit anti-mouse antibody conjugated with horseradish peroxidase was used as the secondary antibody. The blot was subsequently developed with SuperSignal West Femto Western Blot kit (Pierce, Rockford, Ill.). FIG. 2 a shows inhibition of SARS-CoV replication without the pre-treatment of Vero cells; FIG. 2 b shows inhibition effect of the same inhibitors in 2 a with the pre-treatment of Vero cells for 12 h. at 37° C. [0028] FIG. 3 . Vero cells were infected with SARS-CoV and treated with serially diluted concentrations of aurintrucarboxylic acid. HEK293 cells were transfected with an adenovirus construct expressing EGFP. After 24 h., cells were harvested and subjected to 4-12% SDS-PAGE; protein samples were subsequently transferred to PVDF membrane and probed with a mouse monoclonal antibody against SARS-CoV spike protein (for Vero cell extracts), mouse monoclonal antibody against EGFP (for HEK293 cell extracts) and anti-α-actinin antibodies. A rabbit anti-mouse antibody conjugated with horseradish peroxidase was used as the secondary antibody. The blot was subsequently developed with SuperSignal West Femto Western Blot kit (Pierce, Rockford, Ill.). [0029] FIG. 4A . Sequence homology between RDRP, Calpain and Yersina proteins. Frames show putative ATA target sites. [0030] FIG. 4B . Evolutionary associated, sequence similarity of 1QZ0 (1) SARs RSRP (2) and M-Calpain (3). [0031] FIG. 5 . Structural alignment between proteins from SARS-CoV and other RNA viruses. The structural alignment between proteins from SARS-CoV and other RNA viruses based on their 3D atomic coordinate files were performed using Dali and LGA with manual alignment. Solvent inaccessible residues were represented in upper case, while the solvent accessible residues are represented in lower case. The residues that have hydrogen bond to main chain amide are in bold and residues that have hydrogen bond to main chain carbonyl are underlined, Finally, residues that are joined by disulphide bonds are represented in cedilla. 1C2PA—Hepatitis C virus RNA-dependent RNA polymerase; 1DF0A—m-calpain; 1hhsa—bacteriophage phi6 RNA-dependent RNA polymerase; 1KHVA—rabbit hemorrhagic disease virus RNA-dependent RNA polymerase; 1O5SA—SARS-CoV RNA-dependent RNA polymerase; 1QZ0A— Yersinia Pestis phosphatase yoph; 1RDR0—poliovirus RNA-dependent RNA polymerase; 1S4FA—bovine viral diarrhea virus RNA-dependent RNA polymerase; 1SH0A—Norwalk virus RNA-dependent RNA polymerase; 3HVTA—HIV Type 1 RNA-dependent RNA polymerase. [0032] FIG. 6 . Calculated structure (using Autodock) for the interaction of ATA with ypoH. YpoH is a protein tyrosine phosphatase which is essential for virulence in Yersinia pestis. It is known that the functionality of this protein is inhibited strongly by ATA. This figure shows the ten most possible confirmations of ypoh-ATA complexes. The border residues that have contact with ATA are shown and are part of a region that is structurally conserved between ypoh, m-captain and SARS-CoV RdRps, and are constituted mainly of anti-parallel β-strand-turn-β-strand hairpin structures. The two dimensional chemical structure for ATA (C 22 H 14 O 9 ) is shown below. [0033] FIG. 7 . Calculated structure (using Autodock) for the interaction of ATA with m-calpain. The neutral protease (calpain) is a class of cytosolic enzyme that is activated during apoptosis. It is known to be inhibited strongly by ATA. This figure shows the ten most possible confirmations of m-caplain-ATA complexes. The border residues that have contact with ATA are shown and are in a region that is structurally conserved between ypoh, m-caplain and SARS-CoV RdRps and are constituted mainly of anti-parallel β-strand-turn-β-strand hairpin structures. [0034] FIG. 8 . Calculated structure (using Autodock) for the interaction of ATA with RdRp of SARS-CoV. [0035] FIG. 9 . Structure of ATA. [0036] FIG. 10 is a graph showing dynamics of West Nile Virus Replication. [0037] FIG. 11 is a graph showing inhibition of West Nile Virus replication by ATA at increasing concentrations of ATA. [0038] FIG. 12 . Inhibition of vaccinia virus replication by ATA. Hela cells infected with vv (WR) and treated with dilutions of ATA. Plaque counts were done in triplicate on BSC-1 cells, using a 200 μl inoculum and checked at 24 hours post-infection. [0039] FIG. 13 . Time-course inhibition of vaccinia virus replication by ATA. RK13 cells were infected with vv(WR) and the virus titre was checked at various time points. Plaque counts were done in triplicate on BSC-1 cells, using a 200 ul inoculum and checked at 24 hours post-infection. [0040] FIG. 14 . Silver stained gel of the 46 Kdal GST-H1L fusion protein purification. Lane 1: protein mass markers, lane 2: GST column flow through, Lane 3: GST fraction A7 Lane 4: QFF fraction A5, Lane 5: QFF fraction A7, Lane 6: QFF fraction A9. [0041] FIG. 15 . Effect of ATA on H1L catalyzed hydrolysis of pNPP. The reactions were performed at 37 C at 0.06, 0.125, 0.25, 0.5, 1, 2 and 4 mM pNPP. All reactions were performed in triplicate. [0042] FIG. 16 . IC50 analysis of H1L inhibition. Reactions were performed at a substrate concentration equal to the Km value (1 mM), at 37 C. All reactions were performed in triplicate. [0043] FIG. 17 . Comparative inhibition of phosphatase activity by ATA. Activities of each enzyme were standardized, with the same amount of activity being used in the reactions. All reactions were done at 37 C, in serial dilutions of ATA. [0044] Table 1. Vero cells were seeded in a 96-well plate. Dilutions of ATA, interferons α and β were added. After 24 h. 50 μl of Reaction Solution from XTT kit was added to each well and incubated at 37 C for 4 h. Activities of cell proliferations were reflected by spectrophotometric readings. The concentrations of each reagent that inhibits 50% cell proliferation activities (CC 50 ) were compared with the concentrations that inhibit 50% of SARS-CoV replication (EC 50 ), and designated as the selection index (SI). [0045] Table 2. Estimated free energy (final intermolecular energy+torsional free energy) of different targeted proteins in complex with ATA using Autodock. DESCRIPTION OF THE PREFERRED EMBODIMENTS [0046] Unless defined otherwise, all technical and scientific terms used herein have the same meaning as commonly understood by one of ordinary skill in the art to which the invention belongs. Although any methods and materials similar or equivalent to those described herein can be used in the practice or testing of the present invention, the preferred methods and materials are now described. All publications mentioned hereunder are incorporated herein by reference. [0047] As used herein, “effective amount” refers to an amount that is sufficient to achieve the desired result. In regards ATA, an effective amount is capable of inhibiting the target organism. [0048] As used herein, “inhibition” in all its grammatical forms does not imply a complete cessation but rather indicates the activity being inhibited occurs at a lower rate or efficiency. [0049] Described herein is a method of inhibiting replication of an organism which has an essential enzyme which includes a binding groove that is bound by ATA comprising administering to a patient in need of such treatment an effective amount of aurintricarboxylic acid (ATA). [0050] As can be seen in FIG. 4 , a putative 3D structure of SARS RNAP and ATA (using an ATA ligand model) shows the template binding groove of RNAP which is also bound by ATA. That is, ATA binding at this groove prevents template binding and therefore viral replication by RNAP. [0051] Furthermore, as can be seen from the sequence comparison shown in FIG. 5 , enzymes from other organisms known to be inhibited by ATA have similar grooves. It is important to note however that as discussed above ATA has many putative functions and it is possible that it acts through a number of different mechanisms in, addition to the groove binding activity discussed above. [0052] As discussed below, ATA is an effective anti-viral for coronaviruses, for example, SARS-CoV, West Nile Virus, Norwalk, Dengue and Japanese Encephalitis virus. However, we have tried ATA on Ebola and Adenovirus and saw no effect. [0053] In addition, the inhibitory effect of aurintricarboxylic acid (ATA) on vaccinia virus replication in tissue culture is described herein. Concentrations of ATA in the range of 400 ug/mL decreased viral replication as much as 250,000 fold as compared to controls. A block in replication was evident at drug concentrations as low as 25 μg/mL. Inhibition of the viral phosphatase, H1L, which is essential for virus replication, was found to be one mechanism through which ATA exerts its anti-viral effect. The IC50 value of ATA against H1L was found to be 16 μM. This block in enzyme activity leads to a global shutdown in viral protein production and subsequent viral replication. Western blotting also revealed that the ERK signaling cascade was down-regulated in cells treated with ATA. The activity of the ERK signaling cascade has previously been implicated in the vaccinia virus lifecycle. As discussed above, H1L is highly conserved among related viruses and as such ATA is also an effective treatment for infections caused by variola virus, ectromelia virus, monkeypox virus, cowpox virus, myxoma virus and shope fibroma virus. [0054] As discussed above, ATA is a poly-aromatic carboxylic acid derivative, displaying a wide array of biological activities. It has been shown to inhibit endonucleases, topoisomersaes, kinases and phosphatases, as well as apoptosis (26, 48). It has the ability to stimulate the activity of the Insulin-like Growth Factor 1 and Jak2-Stat5 signaling pathways (53). Small molecules like ATA have been shown to interact with receptors at the cell surface (55), where they may alter the activity of signaling pathways. Previous studies have shown that ATA inhibits PTP activity (32, 55). [0055] As discussed herein, ATA inhibited the replication of SARS-Cov. Specifically, when used at concentrations of 800 μg/mL, a 1000 fold decrease in replication was observed, while 400 μg/mL dropped replication 100 fold. Furthermore, when used at 400 μg/mL, ATA could inhibit vaccinia virus replication 250,000 fold, dropping the level of recoverable virus to near non-detectable levels. ( FIG. 13 ). This inhibition was seen in all cell lines tested, including Hela, RK13, Huh7 and vero cells. The virus titre peaked at close to 48 hours post-infection in the untreated controls. This time point also showed the greatest difference in virus titre between treated and untreated cells. Inhibition was seen still seen at ATA concentrations as low as 25 μg/mL. To determine whether ATA has a general antiviral effect, it was tested against Adenovirus. In this case, the drug actually stimulated viral replication. Viral replication was seven times higher when ATA was present compared to no ATA. [0056] The vaccinia virus dual-specific PTP is essential to virus viability. Experiments have shown that vaccina virus mutants lacking a functional copy of H1L have reduced levels of transcription, leading to a resulting lack of replication and infectivity (47). RT-PCR experiments were used to detect this block in viral transcription. The presence of F10 mRNA was decreased in cells treated with ATA. This is consistent with a block in viral transcription. As a control, the mRNA level of GADPH was also followed. This is a host cell house-keeping gene, and a commonly used control. GADPH levels increased in infected cells treated with ATA. This is as expected because during a productive infection, the virus shuts down host protein synthesis. If ATA were inhibiting the virus, you would expect GADPH levels to increase compared to non-treated (ATA) cells. [0057] The homologue of H1L found in other poxviruses has also been shown to be essential (38, 51). H1L has also been implicated in subverting host defense mechanisms in order to allow for a productive infection. Studies have demonstrated that infection of cells with vaccinia virus leads to a reduction in gamma interferon signal transduction. This was found to be a direct result of lower levels of phosphorylated Stat1 in vv infected cells. The phosphorylation of Stat1 is required for its signal transducing activity. Levels of H1L present in the cells correlated inversely with the levels of phosphorylated Stat1. This provides strong evidence that Stat1 is a bona fide substrate for H1L in vivo (52). The gamma interferon signal transduction cascade induces the expression of genes involved in the anti-viral response. This includes the expression of major histocompatibilty complex (MHC) on antigen presenting cells and the activation of natural killer cells (43). The dephosphorylation of Stat1 by H1L may represent one mechanism by which vaccinia virus evades host immune responses. [0058] H1L has also been implicated in other processes that are important for host cell function. Alternative splicing of RNA represents a means by which cells produce diverse and structurally distinct sets of proteins. The SR proteins are responsible for regulating this activity, and phosphorylation in the RS domain has been shown to be essential for proper functioning. When extracts of vaccina virus infected cells were blotted with antibodies against SR proteins, it was found that they were hypo-phosphorylated. Further experiments involved incubating purified SR proteins with purified H1L in an in vitro dephosphorylation assay. This resulted in dephosphorylation of the SR proteins in the RS domain, leading to their inactivation as splicing factors (41). [0059] Several vaccina virus proteins are also substrates for H1L. The A17 and A14 phospho-proteins of vaccinia virus are implicated as being essential viral membrane components. The phosphorylation state of these proteins is regulated by the F10 kinase and H1L. Repression of both A14 and A17 leads to aberrant morphogenesis. However, the role of the phosphorylation of these proteins remains unclear, although it may regulate binding interactions and glycosylation levels (35, 49). [0060] This report is the first to describe kinetic parameters for H1L ( FIG. 15 ). These values are specific to the reaction conditions and the substrate used. Thus, when reactions are performed at a different temperature, pH, or enzyme concentration, the result will vary. Also, the use of a different substrate would produce different results because the affinity of the phosphatase is different for any substrate. Two very important kinetic values for any enzyme are the Vmax and Km values. The Vmax is defined as the maximum rate of an enzyme catalyzed reaction when the active site is saturated with substrate (45). The Vmax for H1L was determined to be 6.40 (abs. units). This value is a reflection of the absorbance reading, and thus, the reaction rate. The Km is defined as the substrate concentration at which half maximal velocity is reached (45). The Km value for H1L was found to be 1 mM pNPP. This means that when pNPP is used at a concentration of 1 mM, the reaction rate will be 50% of the maximum value. The Vmax of H1L was reached at substrate concentrations of 4 mM. These values make allow for determination of the potency of ATA towards H1L. [0061] Phosphatase activity assays were employed to determine if ATA could inhibit H1L. It was found that ATA could block H1L dependant dephosphorylation of pNPP, a commonly used experimental substrate. The IC50 value for this inhibition was 16 μM ( FIG. 16 ). This means that at an ATA concentration of 16 μM, the activity of H1L is reduced by 50%. The control reactions using GST-A30L rule out the possibility of a contaminant phosphatase in the preparation. It also shows that the GST tag does not confer phosphatase activity. ATA had a greater inhibitory effect on the T-cell phosphatase than on H1L, while LAR was less affected. This would mean that using ATA at a concentration effective against H1L would also affect the activity of T-cell PTP in host cells. These assays only detect phosphatase activity towards tyrosine residues. Further experiments need to be done to determine if ATA can also inhibit the Ser/Thr activity of H1L. It should also be noted that the IC50 value against ATA is much lower than the concentration required to inhibit 50% growth in tissue culture. It may be that not all the drug enters the cell, thus decreasing its intracellular concentration. Also, ATA is known to bind to proteins and DNA, decreasing the amount of drug available to inhibit the phosphatase. [0062] Inhibition of H1L will result in a defective infection. Virus transcription will be blocked, to a level proportional to the inhibition of H1L. This means that there will be a reduction in the production of all viral proteins, including H1L itself. This would obviously have a detrimental effect on the virus as proteins necessary for replication, transcription, viral morphogenesis and immune evasion would be found at significantly reduced levels. Host cell gamma-interferon signaling would be restored to wild-type activity, leading to an anti-viral response. H1L inhibition also leads to the hyper-phosphorylation of several viral and host proteins. Thus, the regulatory function of the phosphate groups on these proteins would be lost. [0063] A further aspect of this project was to determine the effect, if any, that ATA exerts on host cell signaling pathways. Vaccinia virus encodes two kinases, B1R and F10L, as well as a phosphatase, H1L. Phosphorylation events are key mediators of host cell signaling cascades, and it seems likely that vaccinia virus exploits these cascades for its own good. Viral interference with host cell signaling is well documented (36, 39, 52, 54). The MAPK (mitogen activated protein kinase) cascade is responsible for controlling cell growth and proliferation, and responses to stress. One route of this cascade, controlled by the extra-cellular regulated kinase 1/2 (Erk1/2), has been shown to be essential for a productive vv infection (57). [0064] Western blotting experiments using an anti-phospho Erk1/2 antibody showed that ATA treatment of cells lead to a decrease in the levels of phospho-Erk protein. Similar results were obtained for the JNK and p38 MAPKs. This phosphorylation is necessary for further transduction throughout the pathway. Thus, this may represent another mechanism by which ATA inhibits vaccinia virus. It is not known, however, at which point in the signaling cascade this block appears at. Inhibition of any kinases upstream of Erk would lead to a similar result. Interestingly, VHR, a human phosphatase related to H1L, as been shown to dephosphorylate Erk1/2 (24, 56). It is possible that H1L may also perform this function, however it would not be expected to lead to complete Erk1/2 dephosphorylation in vivo, as this would block viral infection. It is likely that a balance exists in the phosphorylation state of signaling proteins. Vaccinia virus likely alters this balance, by way of its kinases and phosphatase, to create an intracellular environment more amenable to its own replication. [0065] This study provides evidence that ATA is of therapeutic use against Poxviruses. We have shown that ATA is able to inhibit the phosphatase activity of H1L, with an IC50 value of 16 μM. This enzyme is conserved throughout the Poxviruses, and thus they are all likely to be susceptible to this drug. This enzyme plays essential roles in many aspects of the viral lifecycle, including viral transcription and immune evasion. This makes ATA an ideal drug candidate, as resistance to the drug would difficult to gain. Finally, ATA is able to down-regulate Erk mediated signaling, which is required for vv infection. [0066] As will be appreciated by one of skill in the art, other possible ATA targets can be identified by several different means. For example, the knowledge of the consensus sequence for an ATA binding domain means that sequence and/or structure databases can be searched for enzymes with a similar binding groove. Thus, one aspect of the invention is directed to a method of identifying ATA-target organisms comprising screening a suitable database for matches to a consensus sequence for ATA binding domain, thereby identifying an ATA-inhibitable organism. Another aspect of the invention is directed to organisms identified by this search. [0067] Specifically, protein structural studies can be carried out to investigate the potential binding modes/sites of ATA onto RNA-dependent-RNA-polymerases (RdRp) from SARS-CoV and other pathogenic positive-strand RNA viruses, as well as other proteins in SARS-CoV based on the fact that ATA binds to Ca 2+ -activated neural protease (m-calpain), the protein tyrosine phosphatase (PTP) and HIV integrase which have existing crystal structures. Eight regions with homologous 3D-confirmation were derived from ten proteins of interest. One of the regions, R binding (754-766 in SARS-CoV's RdRp), located in the palm sub-domain consists of mainly anti-parallel β-strand-turn-β-strand hairpin structures that include two of the three RdRp catalytic sites (Asp 706, Asp 761), was also predicted by a molecular docking method (based on free energy binding of ΔG) to be an important binding motif recognized by ATA. The existence of this strictly conserved region that includes catalytic residues, coupled with the homologous ATA binding pockets and their consistent ΔG values indicates that ATA is involved in an analogous inhibition mechanism in SARS-CoV's RdRp, m-calpain, PTP and HIV integrase. Furthermore, as discussed below, other ATA-inhibited organisms can be identified using the searching method described herein. In another aspect, the invention is directed to this method of identifying organisms having similar motifs in key enzymes as well as the use of ATA to inhibit growth of the organisms so identified. [0068] Alternatively, ATA-sensitive organisms can be determined by growing an organism under normally reproductive conditions in the presence of ATA and then determining if reproduction has occurred. This may be done by comparison with an untreated control. For example, if the organism is a bacteria or fungus, the organism may be grown in culture media at an appropriate incubation temperature. If the organism of interest is a virus, the organism is grown under appropriate culture conditions in the presence of ATA, virus is harvested and viral numbers are determined, for example, by a plaque assay. [0069] According to another aspect of the invention, there is provided a method of inhibiting replication of an organism having an essential enzyme having an ATA binding domain comprising administering to that organism an effective amount of ATA. [0070] As used herein, “aurintricarboxylic acid derivatives” or “ATA derivatives” includes chemically modified variants thereof, that is, compounds having the same general structure as ATA and further having the same biological activity as ATA and pharmaceutically acceptable salts thereof, wherein “biological activity” refers to the inhibition of viral replication. [0071] In a preferred embodiment, ATA is used as a pharmaceutical composition for treating a patient suffering from or at risk of developing or suspected of having a coronavirus infection. Administration of an effective amount of ATA will have at least one of the following effects: inhibition of RNA synthesis, reduction of viral load, and amelioration of associated symptoms. [0072] As will be appreciated by one of skill in the art, members of the coronavirus family include but are by no means limited to Bovine viral diarrhoea virus, feline infectious peritonitis, Porcine epidemic diarrhea virus, mouse hepatitis virus, human coronavirus 229E and human coronavirus OC43. [0073] In another embodiment, ATA is used as a pharmaceutical composition for treating a patient suffering from or at risk of developing or suspected of having a West Nile Virus infection, a Norwalk virus infection, a Dengue virus infection, or a Japanese Encephalitis infection. Administration of an effective amount of ATA will have at least one of the following effects: inhibition of RNA synthesis, reduction of viral load, and amelioration of associated symptoms. Specifically, ATA has been shown to inhibit West Nile Virus 10000 fold as can be seen in FIGS. 10 and 11 . [0074] In a preferred embodiment, ATA is used as a pharmaceutical composition for treating a patient suffering from or at risk of developing or suspected of having a poxvirus infection. Administration of an effective amount of ATA will have at least one of the following effects: inhibition of RNA synthesis, reduction of viral load, and amelioration of associated symptoms. As discussed above, members of the poxvirus family include but are by no means limited to variola virus, ectromelia virus, monkey pox virus, cowpox virus, myxoma virus and shope fibroma virus. [0075] It is of note that ATA discussed above may be prepared to be administered in a variety of ways, for example, topically, orally, intravenously, intramuscularly, subcutaneously, intraperitoneally, intranasally or by local or systemic intravascular infusion using means known in the art and as discussed below. [0076] It is of note that as discussed herein, the ATA may be arranged to be delivered at a dosage of about 25 mg to about 800 mg per kg of the subject, or about 50 mg to about 800 mg per kg of the subject, or about 100 mg to about 800 mg per kg of the subject or about 200 mg to about 800 mg per kg of the subject or about 400 mg to about 800 mg per kg of the subject. As will be apparent to one knowledgeable in the art, the total dosage will vary according to the weight of the individual as well as other factors, for example, the age and condition of the patient. As discussed above, in some embodiments, an effective amount of ATA may be an effective concentration of between 10-800 μg/ml, between 50-800 μg/ml, between 100-800 μg/ml or between 400-800 μg/ml. [0077] In some embodiments, ATA at concentrations or dosages discussed above may be combined with a pharmaceutically or pharmacologically acceptable carrier, excipient or diluent, either biodegradable or non-biodegradable. Exemplary examples of carriers include, but are by no means limited to, for example, poly(ethylene-vinyl acetate), copolymers of lactic acid and glycolic acid, poly(lactic acid), gelatin, collagen matrices, polysaccharides, poly(D,L lactide), poly(malic acid), poly(caprolactone), celluloses, albumin, starch, casein, dextran, polyesters, ethanol, mathacrylate, polyurethane, polyethylene, vinyl polymers, glycols, mixtures thereof and the like. Standard excipients include gelatin, casein, lecithin, gum acacia, cholesterol, tragacanth, stearic acid, benzalkonium chloride, calcium stearate, glyceryl monostearate, cetostearyl alcohol, cetomacrogol emulsifying wax, sorbitan esters, polyoxyethylene alkyl ethers, polyoxyethylene castor oil derivatives, polyoxyethylene sorbitan fatty acid esters, polyethylene glycols, polyoxyethylene stearates, colloidol silicon dioxide, phosphates, sodium dodecylsulfate, carboxymethylcellulose calcium, carboxymethylcellulose sodium, methylcellulose, hydroxyethylcellulose, hydroxypropylcellulose, hydroxypropylmethycellulose phthalate, noncrystalline cellulose, magnesium aluminum silicate, triethanolamine, polyvinyl alcohol, polyvinylpyrrolidone, sugars and starches. See, for example, Remington: The Science and Practice of Pharmacy, 1995, Gennaro ed. [0078] As will be apparent to one knowledgeable in the art, specific carriers and carrier combinations known in the art may be selected based on their properties and release characteristics in view of the intended use. Specifically, the carrier may be pH-sensitive, thermo-sensitive, thermo-gelling, arranged for sustained release or a quick burst. In some embodiments, carriers of different classes may be used in combination for multiple effects, for example, a quick burst followed by sustained release. [0079] The invention will now be described by examples; however, the invention is not limited to the examples. Cell Culture and Viral Plaque Assay [0080] The African green monkey kidney cell line Vero cells were cultured in Dulbecco's Modified Eagle's Medium, supplemented with 10% heat-inactivated fetal bovine serum (Inivitrogen, Carlsbad Calif.), 1% penicillin/streptomycin and 10 mM HEPES (pH=7.2). Vero cells have been shown to be susceptible to SARS-CoV infection [4]. All cell cultures were maintained in a humidified 5% CO 2 incubator at 37° C. In all experiments, the multiplicity of infection (MOI) was 0.01. Aurintrycarboxylic acid (Sigma, St. Louis, Mo.) was prepared in the culture media and added into the samples in serial dilutions comprised of 0.8 mg/ml, 0.4 mg/ml, 0.2 mg/ml 0.1 mg/ml and 0.05 mg/ml. Plaque assays were performed 24 h. post infection using procedures as described [16]. We also performed inhibition analysis with interferon α and β both at the concentration of 5000 IU/ml [10,12]. Western Blotting [0081] Protein samples from Vero cell extracts were fractionated on 4-12% SDS-PAGE (Invitrogen Carlsbad Calif.) and transferred to PVDF membrane using a semi-dry protein transfer apparatus (Bio-Rad, Hercules Calif.). The membrane was blocked for 1 hr with 5% skim milk in TBS buffer (20 mM Tris base, 137 mM NaCl, pH7.6) containing 0.2% Tween-20. The membrane was then probed with a mouse monoclonal antibody against SARS-CoV spike protein. Rabbit anti-mouse HRP-conjugated antibody (Amershambiosciences, Piscataway, N.J.) was subsequently added for an additional incubation of 1 hr at 37° C. The results were revealed using Pierce Biotechonlogy SuperSignal West Femto Maximum Sensitivity Substrate (Rockford, Ill.). Real-Time RT-PCR Analysis [0082] The analysis was performed with Prism 7700 real-time PCR instrument from PerkinElmer (Wellesley, Mass.) following the manufacturer's protocol. Supernatant samples from SARS-CoV infected Vero cells were collected for viral RNA extractions using RNeasy kit from Qiagen (Valencia, Calif.). Primers and probe used in the RT-PCR are as follows: Probe: 6FAM-ACCCCAAGG TTTACCC (SEQ ID No. 1); Forward: ACCAGAATGGAGGACGCAATG (SEQ ID No. 2); and Reverse: GCTGTGAAC CAAGAC GCAGTATTA T (SEQ ID No. 3). Inhibition of Adenovirus Expressing EGFP [0083] ATA was used in a comparative study for inhibition of the adenovirus Adeno X, which carries a replication reporter gene expressing EGFP (Clontech, Palo Alto, Calif.). Approximately 5 mg of adenovirus X were transfected into HEK293 cells, using Effectene (Qiangen, Valencia, Calif.), followed by addition of serial dilutions of ATA-Western blot analysis was followed using the antibody against EGFP (Clontech, Palo Alto, Calif.). Cell Proliferation Assay [0084] The XTT kit (Roche, Mannheim, Germany) was used to measure the toxicity of ATA. Briefly, Vero cells were seeded in a 96-well plate, dilutions of ATA and interferons were added to the cells and incubated for 24 h. The colorimetric detection reagent from the XTT kit was subsequently added to the cells. Results were determined by spectrophotometer at 450 nm. Inhibition of SARS-CoV Replication by ATA [0085] We have tested the anti-SARS-CoV effect of aurintricarboxylic acid (ATA) over a wide range of concentrations, i.e. 0.8 mg/ml, 0.4 mg/ml, 0.2 mg/ml, 0.1 mg/ml and 0.05 mg/ml prepared in Minimum Essential Medium (MEM) with 10 mM HEPES (pH 7.2) and 10% of fetal calf serum. Plaque assays were used to determine the effect of ATA on SARS-CoV replication. Vero cells were infected with SARS-CoV in a 24-well plate; serial dilutions of ATA were added to the infected cells after the initial virus adsorption step. Twenty-four hours post infection, we collected the supernatants from the aforementioned cultures for plaque assays to determine the inhibitory effect of ATA on SARS-CoV replication. As shown in FIG. 1 , in comparison with the untreated cells, more than 1,000-fold inhibition of virus replication was observed when the culture was treated with ATA at a concentration of 0.8 mg/ml, while at least 100-fold inhibition was observed at 0.4 mg/ml. An inhibitory effect on viral replication could still be observed at 0.2 mg/ml, with viral replication level being 10 times lower than that of the control. [0086] We also compared the inhibitory effect of ATA on SARS-CoV replication with that of IFN α and β. To this end, dilutions of ATA and IFNs at the highest effective concentrations were used to treat the cultures [10,12], followed by determination of virus loads. The virus loads at 24 h post-infection were quantified by real-time RT-PCR analysis using specific primers and probes against SARS-CoV nucleocapsid protein. As shown in FIG. 2 a , ATA at 0.8 mg/ml inhibited the virus RNA replication by more than 1,000 fold, versus 100 fold inhibition by interferon α at 5000 IU/ml and 10 fold by interferon β at 5000 IU/ml ( FIG. 2 a ). This result indicates that ATA was about 10 times more potent than interferon α and 100 times more potent that interferon β as an anti-SARS-CoV agent. [0087] To further analyze whether there is a prophylactic effect of ATA, we pre-treated the cells with a series of concentrations of ATA, interferons α and β for 12 h before the adsorption of SARS-CoV to Vero cells and then added the above inhibitors after adsorption. As shown in FIG. 2 b , the inhibition effect of ATA and interferons α and β were about the same as samples without the pre-treatment, implicating that the inhibition effect may take place after the virus enter the cells. Western Blot Analysis [0088] To further confirm the inhibitory effect, we performed Western blot analysis using a monoclonal antibody against SARS-CoV spike (S) protein. As shown in FIG. 3 a , the level of the S protein was significantly lower in the ATA treated group than in the untreated group. At a concentration of 0.8 mg/L, ATA blocked viral protein synthesis, confirming that ATA can significantly inhibit viral protein synthesis. [0089] To characterize the specificity of the anti-viral activity of ATA, we also tested the compound for its ability to block protein expression by adenovirus replication. The replication deficient adenovirus type 5 expressing EGFP was used to infect HEK-293 cells. The same concentrations of ATA used in the above-mentioned SARS-CoV inhibition experiments were added to the adenovirus-infected cells. A Western blot analysis was subsequently performed to determine the expression level of EGFP. No significant inhibition was observed in any ATA-treated samples compared with non-ATA-treated cells ( FIG. 3 b ), indicating that the inhibition of SARS-CoV replication by ATA was clearly selective. Cell Proliferation Assay [0090] To further explore the therapeutic potentials of ATA, we determined the selectivity index (SI) as defined by the ratio of drug concentration causing cellular toxicity to that producing anti-viral effect. To this end, non-radioisotope cell proliferation analysis system XTT from Roche (Mannheim, Germany) was used. CC 50 indicates the concentration that causes 50% of the cytotoxicity, while EC 50 means the concentration of inhibitors that inhibited 50% of the virus replication. We found that the SI of ATA is 187 versus 30 of IFN α and 20 of IFN β; indicating ATA is a potent anti-viral compound with low toxicity (table 1). [0091] Because of its low toxicity in cell culture and animals [17, 18], ATA has been evaluated for its anti-viral activities in viruses such as immunodeficiency virus type I [17, 19]. The potency of ATA against SARS-CoV replication appears to be higher that that of the reported chemicals such as glycyrrhizin and recently reported nelfinavir [12, 20]; both drug candidates reported to have two logs or less inhibition effect on SARS-CoV replication, while ATA showed more than three logs of inhibition effect. The biological activities of ATA are believed to be quite complicated, including inhibition of protein synthesis, prevention of the attachment of mRNA to ribosomes in cell-free systems and suppression of enzymes involved in polynucleotide metabolism [21]. [0092] As will be apparent to one of skill in the art, vaccine development could take years to complete and serious adverse reactions have been reported in other coronavirus vaccine studies, i.e., exacerbation of disease in animals receiving vaccines prior to infection [22]. Certain precautions have been proposed for the development of SARS-CoV vaccines due to potential detrimental effects [23], meaning that the search for anti-SARS-CoV drugs must be pursued. [0093] The 3D theoretical model for RdRp (RDRP, ID=105S, Xu et al., 2003, NAR 31: 7117-7130), spike protein subunit 1 (S 1 , ID=1Q4Z, Spiga et al., 2003, Biochem Biophys Res Commun 310: 78-83), spike protein subunit 2 (S 2 , ID=1Q4Y, Spiga et al., 2003) of SARS-CoV, the 3D crystal structure for nucleocapsid protein (N protein, ID=1SSK, Huang et al., 2004, Biochemistry 43: 6059-6063), non-structural protein 9 (Nsp9, ID=1QZ8, Egloff et al., 2004, PNAS 101: 3792-3796), main protease (3Clpro, ID=1Q2W, Anand et al., 2003, Science 300: 1763-1767) of SARS-CoV, the crystal structure for YopH from Yersinia pestis (ID=1QZ0, Sun et al., 2003, J Biol Chem 278: 33392-33399), m-calpain from Rattus norvigecus (ID=1DF0, Strobl et al., 2000, PNAS 97: 588-592), RNA-dependent RNA polymerase from Dsrna Bacteriophage φ6 (ID=1HI8), Rabbit Hemorrhagic Disease Virus (ID=1KHW), Poliovirus (ID=1RDR), Bovine Viral Diarrhea Virus (IVDV, ID=1S4F), Norwalk virus (ID=1SH0) and HIV reverse transcriptase (ID=3HVT) were downloaded from Protein Data Bank. PRODRG2 (Aalten, 2004, Acta Crystallographica D60, in press) and JME editor (http://www.cem.msu.edu/˜reusch/VirtualText/Questions/MOLEDITOR/jme_window.html) were used to derive the atomic coordinates of ATA. VAST and DALI programs were used to locate similar structural patterns between crystal structures of yopH and RdRps. Sequential structural alignment was done by CE (Shindyalov and Bourne, 1998, Protein Eng 11: 739-747) and COMPARER (Sali and Blundell, 1990, J Mol Biol 212: 403-428). Finally, 3D structural comparative analysis was performed by LGA (Zemla, 2003, NAR 31: 3370-3374). Preparation of macromolecule and ligand prior molecular docking was done using WhatiF software (Vriend, 1990, J Mol Graph 8:52-56). Molecular docking to determine the best confirmation in terms of lowest Gibbs free energy and shape complementarity was performed using Autodock 3.0 (Morris et al., 1998, Journal of Computational Chemistry 19: 1639-1662). The visualization of the 3D structural data was generated by Rasmol (Bernstein, TIBS 25: 453-455). [0094] Protein 3D structural alignments were performed on ten amino acid (mainly RdRps) sequences of interest, namely: Hepatitis C virus RNA-dependent RNA polymerase; m-caplain; bacteriophage φ6 RNA-dependent RNA polymerase; rabbit hemorrhagic disease virus RNA-dependent RNA polymerase; SARS-CoV RNA-dependent RNA polymerase; Yersinia pestis phosphatase yopH; poliovirus RNA-dependent RNA polymerase; bovine viral diarrhea virus RNA-dependent RNA polymerase; Norwalk virus RNA-dependent RNA polymerase; and HIV Type 1 RNA-dependent RNA polymerase. Regions with homologous 3D-conformations were identified together with their conserved secondary structures (shown in FIG. 5 ). In total, there are eight structurally conserved motif blocks (CMBs) with each block extending at least eight amino acid residues. The secondary structures for all CMBs include six α-helices and two β-strands regions. The exact position of all CMBs in each protein are provided in FIG. 5 . [0095] Analysis of the molecular docking method based on free energy of ligand binding, ΔG, between ATA and all proteins ( FIGS. 6 and 7 ) revealed that ATA binds favorably to one structurally conserved region (R binding ) among all proteins ( FIG. 5 ). As shown in FIG. 8 , the corresponding R binding region in SARS-CoV's polymerase (Ser 754-Tyr 766) overlapped with one CMB. R binding is located in the palm sub-domain and consists mainly of anti-parallel β-strand-turn-β-strand hairpin structures. This conserved region is similar to the majority of the remaining nine proteins in terms of their secondary structures. Surprisingly, this R binding region also contains a highly conserved ‘XSDD’ amino acid motif that is especially prominent among viral RdRps, of which two of the highly conserved aspartic acid (D) residues form the catalytic center important for polymerase activity (Xu et al., 2003). [0096] The free energy of ligand binding (final intermolecular energy+torsional free energy), ΔG, between ATA and all proteins is shown in Table 2. The proteins that were documented to be inhibited by ATA were assigned as positive controls (HIV integrase, yopH and m-calpain), their estimated free energies of binding were −11.88 kcal/mol, −7.79 kcal/mol and −7.67 kcal/mol respectively. Any estimated free energies of binding approximately −7.67 kcal/mol or lower are candidates for inhibition by ATA if the corresponding binding motif includes catalytic sites of that specific protein. When we studied the binding of ATA onto various RNA dependent RNA polymerases from other organisms, the free energy of binding for most RdRps were significantly higher, than −7.67 kcal/mol, and therefore a lower inhibition by ATA (Bovine viral diarrhea virus, Dsrna bacteriophage, Feline calicivirus, Hepatitis C virus, HIV, Poliovirus and Rabbit hemorrhagic disease virus). Only the RdRps from SARS-CoV (ΔG=−7.68) and Norwalk virus (ΔG=−14.92) were estimated to have a lower ΔG (<−7.67 kcal/mol). [0097] It is of note that as discussed above, the R binding domain includes 2 of the 3 predicted RdRp catalytic residues and a highly conserved secondary structure. While not wishing to be bound to a particular theory, the inventors believe that these residues are important for metal ion chelation (Bressanelli et al., 2002, J Virol 76: 3482-3492; Beese and Seitz, 1991, EMBO J 10: 25-33; Huang et al., 1998, Science 282: 1669-1675). Structurally, this binding pocket is located in the palm sub-domain which consists mainly of anti-parallel β-strand-turn-β-strand hairpin structures. [0098] As reported above, there are eight structurally conserved CMBs and their secondary structures include six α-helices and two β-strand regions. Among these structurally conserved regions, we subsequently identified that there was one common region recognized by ATA. The binding of ATA to this region also fulfilled the lowest free energy of ligand binding. The beauty of using this free energy of ligand binding is that we were able to quantify the binding strength between a macromolecule and a ligand. Therefore, if a ligand binds strongly (with lower free energy of ligand binding ΔG) to the active domains of one specific protein, it will presumably inhibit the activity for that specific protein. On the other hand, if ATA does not bind to the active domains of the protein, ATA will not be able to inhibit the function of the protein regardless of the strength of binding. One Step Growth Experiment [0099] Hela cells were seeded in a 6 well plate and infected with vaccinia virus (WR) at a multiplicity of infection of 5. After a one hour incubation at 37 C and 5% CO 2 , the media was removed and the cell monolayer was washed three times with cold phosphate buffered saline. Next, appropriate dilutions of ATA were made in 1 mL of media and added to the cells. At various time intervals, the cells were removed with a cell scraper and frozen at −80 C. A freeze/thaw method was used to extract the virus from the cells. The supernatant from the freeze/thaw was centrifuged at 2000×g for 5 minutes and the supernatant was collected. Titration of the virus was done on BSC-1 cells. A one-in-ten dilution of the supernatant was made, and 200 μl was applied to the cell monolayer. Cells were fixed after a two day incubation at 37 C, in 5% CO 2 . 3%, folmaldehyde was used as a fixative, and the cells were stained with 0.5% crystal violet. The virus titre was calculated by counting the number of plaques, and determining the number of plaque forming units per mL. Protein Expression [0100] The vaccinia H1L and A30L open reading frames were amplified by PCR (polymerase chain reaction) from vaccinia virus WR genomic DNA. The H1L N-terminal primer was 5′-TAAAGGATCCATGTACCCATACGATGTTCCAGATTACGCTATGGATAAGAAAA GTTTGTATAAA-3′ (SEQ ID NO. 4). The H1L C-terminal primer was 5′-TTTATACAATAACTATTCTTAATTGAGCTCGCCT-3′ (SEQ ID NO. 5). The A30L N-terminal primer was 5′-AATTGGATCCATGTACCCATACGATGTTCCAGATTACGCTATGGAAGACCTTA ACGAGGCAAACT-3 (SEQ ID NO. 6); and the A30L C-terminal primer was 5′-CCTTCTCTTAAGTTAGCAGCAACTGAGCTCAAAT-3′ (SEQ ID NO. 7). Primers contain a 5′ BamH1 and 3′ Xho1 restriction sites to facilitate cloning. PCR products were ligated into pCR 2.1-TOPO (Invitrogen) and sequenced for fidelity. Coding regions for the genes were then excised with BamH1 and Xho1 and ligated into pre-digested (BamH1/Xho1) pGEX-6P-3. These plasmids were then used to transform E.coli. BL21-pLys-S-DE3. This strain produces lysozyme, which aides in disrupting the cell membrane for protein purification. 500 mL cultures were inoculated with an overnight culture and allowed to grow at 37 C until an OD of 0.7 was reached at 600 nm. Isopropyl-B-D-thiogalactopyranoside (IPTG) was then added to a final concentration of 1 mM and expression of GST-A30L and GST-H1L fusion proteins was continued for 3 hours. For reasons unknown, a fresh transformation of the H1L-pGEX construct into the cells had to be carried out for each new purification. Plating out a glycerol stock of the recombinant bacteria, and performing protein expression from here resulted in low quantities of protein. Protein Purification [0101] Induced cultures expressing GST-A30L and GST-H1L were pelleted by ultra-centrifugation at 5000×g for 15 minutes. Pellets were then lysed in 50 mM Tris-HCl, pH 8.0, 150 mM NaCl, 2 mM EDTA and 1 mM DTT. Mini-Tab (Roche) protease inhibitor tablets were also added. Extra lysozyme was added to 5 mM in order to facilitate cell lysis. The solution was then sonicated, to shear the DNA present, which can clog purification columns. The soluble protein fraction was isolated by ultra-centrifugation at 18000 rpm for 25 minutes, and filtered through a 0.45 μM syringe filter (Corning). Fractions were then loaded onto an AKTA UPC-900 fast performance liquid chromatography (FPLC) unit (Amersham). Protein samples were run through a GSTrap FF purification column (Amersham) at a flow rate of 0.5 mL/min. The column was previously washed with 0.01M Phosphate buffered saline (PBS), pH 7.2. The GST fusion proteins were eluted in 1 mL fractions with 50 mM Tris-HCl, 10 mM reduced glutathione, pH 8.0. The purity of these preparations was assessed by SDS-polyacylamide gel electrophoresis and determined to be approximately 90%. The fraction containing the protein of interest was then further purified with QFF anion exchange columns (Amersham). The protein was loaded in the GST column elution buffer, and eluted in 50 mM Tris-HCl via a salt gradient. Samples of the eluted fraction were then run on a gel and the fractions containing the protein of interest were stored. Purity was estimated to be greater that 95%. RT-PCR Analysis [0102] mRNA levels were analyzed using an Applied Biosystems 7300/7500 Real Time PCR System, following the manufacturer's protocol. Hela cells were grown in the presence or absence of ATA for 3 hours, and supernatant samples were collected at 0, 2, 6, 12 hrs. RNA was purified with an Rneasy Kit (Qiagen). 1 μg of viral RNA was used to detect F10L mRNA levels. As a control, the levels of GADPH were also followed. H1L Kinetic Analysis [0103] Before characterizing the inhibition of H1L by ATA, it was necessary to determine the Km and Vmax values of the enzyme. These kinetic parameters are then used to optimize conditions for an IC50 inhibition assay. To obtain these values, enzyme activity assays were performed in 96 well plates. The assays were performed in 25 mM HEPES, 50 mM NaCl, 5 mM DTT, 2.5 mM EDTA in a final volume of 200 μL. The first step involved the addition of buffer, and 1.5 μg of enzyme. The Bradford Method was used to quantify the purified protein. Additionally, 25 μg of BSA was added to stabilize the enzyme and create optimal kinetic conditions. Reactions were pre-incubated at 37 C for 15 minutes in a water bath. Substrate was then added to start the reaction. Obtaining the Vmax and Km values requires the plotting of a Michealis-Menton curve, in which then enzyme concentration is held constant while the substrate concentration varies. Substrate concentrations used were 0.06, 0.125, 0.25, 0.5, 1, 2 and 4 mM pNPP. These were the concentration reached in the final 200 μl reaction. A further incubation of ten minutes followed addition of substrate. The reaction was stopped by the addition of 13% K 2 HPO 4 . Absorbances were read with a SpectraMax Plus (Molecular Devices), at 405 nm. The Km and Vmax values were calculated from the Lineweaver-Burk curve of the data. The substrate, pNPP (p-nitrophenylphosphate), absorbs at 405 nm, so a set of control reactions was also run. These involved addition of substrate but no enzyme in the first step. The values obtained in the controls were subtracted from the appropriate samples. IC50 Determination [0104] Assays were performed in the buffer previously mentioned. Reactions were carried out at a fixed substrate concentration, equal to the Km value of the enzyme. ATA was used at varying concentrations, ranging from 2 to 125 μM. The reaction sequence is the same as described above. Since ATA and pNPP absorb at 405 nm, a series of control reactions, representing the dilutions of ATA and pNPP used were set up. These reactions did not receive enzyme, and the absorbance value read from them was then subtracted from the appropriate samples. The absorbance at 405 nm was plotted against the concentration of ATA used. The IC50 value is that concentration at which half-maximal enzyme activity occurs. ATA Inhibits Replication of Vaccinia Virus in Tissue Culture [0105] The effect of ATA on vv replication was studied in numerous cell lines using plaque assays. A range of concentrations from 0 to 400 μg/mL was used. Hela cells were infected with vv (WR) at a MOI of 10 ( FIG. 12 ). The viral load dropped by over one log value when 400 μg/mL of ATA was used compared to the control with no drug. Replication inhibition was seen in the lowest dilutions used. Infectious virus counts steadily decreased as the concentration of the drug was increased. Huh7, Vero and Rk13 cell lines were also used to repeat the same experiments. In each case, the results were similar. All experiments were performed in triplicate. [0106] To further characterize this inhibition, a time-course inhibition assay was performed with RK13 cells. The cells were infected with vv (WR) at a MOI of 5. ATA was applied at various concentrations, and the PFU/mL was determined at 0, 4, 24, 48, 72 and 96 hours post-infection ( FIG. 13 ). In the absence of ATA, the levels of virus increased to a maximum of 1.55×10 pfu/mL at 48 h.p.i., and then subsequently decreased. When ATA was used at 400 ug/mL, the viral load again peaked at 48 h.p.i., with the plaque count being 616 pfu/mL. This represents a 250,000 fold difference in treated versus untreated cells. At the final time point (96 hrs.), the viral count in the treated sample was only 0.00026% of the viral count in the untreated control. ATA Inhibits the H1L Phosphatase [0107] ATA has previously been shown to inhibit phosphatases, including those from both eukaryotic and prokaryotic sources. To further characterize the mechanism behind the replication inhibition, it was decided to perform enzyme inhibition assays. These assays specifically detect phosphatase activity towards tyrosine residues. The VV H1L phosphatase was purified to near homogeneity ( FIG. 14 ) as described in materials and methods. [0108] All assays were performed using the GST-H1L fusion protein. Previous studies have shown that the fusion protein and H1L alone have nearly identical activities (37). As a control, preparations of GST-A30L, were also run in identical phosphatase assays. These reactions always gave results comparable to the blank. Assays were first performed to determine the Vmax and Km of H1L towards the experimental substrate, para-nitrophenylphosphate. The product released, para-nitrophenolate, results in the yellow colour of the reaction (33), which is quantified by following the absorbance at 405 nM. The resulting data collected was fit to a Lineweaver-Burk plot, with the Vmax being 6.40 (abs. units) and a Km of 1.00 mM pNPP ( FIG. 15 ). These values are solved from the formula provided in FIG. 17 , with the Vmax representing the y-axis value and the Km represented by the x-axis value. [0109] To determine the extent to which ATA inhibits the activity of H1L, IC50 assays were performed. The IC50 value was found to be 16 μM ( FIG. 16 ). The IC50 value of ATA against other PTPs (protein tyrosine phosphatases) has been previously described (32, 46). [0110] The potency of ATA towards H1L and other phosphatases was also compared. The enzyme activity of H1L, LAR, YopH and T-cell phosphatases was standardized, such that for the assay, the same amount of enzyme activity was added to each reaction. It was found that ATA was more selective for the YopH and T-cell phosphatase and less so for LAR, when compared to H1L ( FIG. 17 ). Activity of the YopH and T-cell phosphatases only became detectable at the lowest concentration of ATA used, 4 μM, while H1L and LAR still showed activity at the highest concentration of 124 μM. [0111] While the preferred embodiments of the invention have been described above, it will be recognized and understood that various modifications may be made therein, and the appended claims are intended to cover all such modifications which may fall within the spirit and scope of the invention. REFERENCES [0000] [1] Marra et al., 2003, Science. 300 (5624), 1399-1404. [2] Rota et al., 2003. Science. 300 (5624), 1394-1399. [3] Nicholls et al., 2003, Respirology. Suppl:S6-8. [4] Holmes K V,. 2001 Coronaviruses. In: Knipe D M and Howley P M (Eds), Field's Virology. Lippincott, Williams & Wilkins, Philadelphia, Vol. I, pp. 1187-1203. [5] Hierholzer, 1976, Virology. 75(1):255-65. [6] Bucknall et al., 1972, Proc Soc Exp Biol Med. 139(3):811-7. [7] Van Der Hoek et al., 2004, Nat Med. 10(4):368-73. [8] He et al., 2003, Biochem Biophys Res Commun. 28;311(4):870-6. [9] He et al., 2004, Biochem Biophys Res Commun 316(2) 476-483 [10] Tan et al., 2004, Emerg Infect Dis. 10(3) 581-586. 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[27] Broyles, 2003, J. Gen. Virol. 84, 2293-2303. [28] Broyles et al., 2004, J. Virol. 78, 2137-2141. [29] Byrd et al., 2004, J. Virol. 78, 12147-12156. [30] Clercq and Neyts, 2004, Rev. Med. Virol. 14, 289-300. [31] Clercq, 2001, Clin. Mic. Rev. 14, 382-397. [32] Cho et al., 2004, Mol. Cells. 18, 46-52. [33] Denu et al., 1995, J Biochem. 270, 3796-3803. [34] Denu and Dixon, 1995, Proc. Natl Acad. Sci. 92, 5910-5914. [35] Derrien et al., 1999, J. Virolgy. 73, 7287-7296. [36] Diperna et al., 2004, J. Bio. Chem. 279, 36570-36578 [37] Guan et al., 1991, Nature. 350, 359-362. [38] Hakes et al., 1993, Proc. Natl. Acad. Sci. 90 4017-4021. [39] He et al., 2003, Biochem. Biophy. Res. Comm. 311, 870-876. [40] He et al., 2004, Biochem. Biophys. Res. Comm. 320, 1199-1203. [41] Huang et al., 2002, EMBO rep. 3, 1088-1093. [42] Hunt and Wagner, 1975, J. Virol. 16, 1146-1153. [43] Karp, 2002, Cell and molecular biology. John Willey, New York. [44] Kern, 2003, Antiviral Reas. 57, 35-40. [45] Lehninger A., Nelson D., Cox. (1993). Principles of Biochemistry. Worth publishers, New York, N.Y. [46] Liang et al., 2003, J. Biochem. 278, 41734-41741. [47] Liu et al., 1995, J. Virol. 69, 7823-7834. [48] McCune et al., 1989, Biochem. J. 259, 925-927. [49] Mercer and Traktman, 2003, J. Virol. 77, 8857-8871. [50] Moss, 2001, Poxviridae: The viruses and their replication. Fields Virology. 2849-2855. [51] Mossman et al., 1995, Virol. 206, 572-582. [52] Najarro et al., 2001, J. Virol. 75, 3185-3196. [53] Rui et al., 1998, J. Biochem. 273, 28-32. [54] Shisler and Jin, 2004, J. Virol. 78, 3553-3560. [55] Tian et al., 1998, Science. 281, 257-259. [56] Todd et al., 1999, J. Bio. Chem. 274, 13271-13280. [57] Wang et al., 2004, Nature Imm. 1-9 [0000] TABLE 1 Selection of Inhibitors CC50 EC50 SI ATA 37.5 mg/ml 0.2 mg/ml 187 IFNα 15000 IU/ml 500 IU/ml 30 IFNβ 10000 IU/ml 500 IU/ml 20 [0000] TABLE 2 Estimated Free energy (final intermolecular energy + torsional free energy) of different targeted proteins in complex with ATA using Autodock. Estimated Free Energy of Binding in Complex Type of with ATA/ Targeted Protein structure used kcal/mol RNA RdRp (Bovine Viral Crystal Structure −5.82 dependent Diarrhea Virus) RNA RdRp (Dsrna Bacterio- Crystal Structure −5.87 polymerase phage) RdRp (Feline calicivirus) Theoretical model −5.83 RdRp (Hepatitis C Virus) Crystal Structure −5.8 RdRp (HIV) Crystal Structure −5.82 RdRp (Norwalk virus) Crystal Structure −14.92 RdRp (Poliovirus) Crystal Structure −5.83 RdRp (Rabbit Hemor- Crystal Structure −5.94 rhagic Disease Virus) RdRp (SARS-CoV) Theoretical model −7.68 Protein Yoph* Crystal Structure −7.79 known to HIV integrase* Crystal Structure −11.88 be inhibited m-Calpain* Crystal Structure −7.67 by ATA* Other Main Protease (3CL) Crystal Structure −7.78 SARS' N Protein Crystal Structure −8.32 proteins Nsp 9 Crystal Structure −7.64 S 1 [28] Theoretical model −18.59 S 2 [28] Theoretical model −7.66 S 1 (unpublished model) Theoretical model −14.79 S 2 (unpublished model) Theoretical model −15.22

The severe acute respiratory syndrome virus (SARS) is a coronavirus that instigated regional epidemics in Canada and several Asian countries in 2003. The newly identified SARS coronavirus (SARS-CoV) can be transmitted among humans and cause severe or even fatal illnesses. As preventive vaccine development takes years to complete and adverse reactions have been reported to some veterinary coronaviral vaccines, anti-viral compounds must be relentlessly pursued. In this study, we analyzed the effect of aurintricarboxylic acid (ATA) on SARS-CoV replication in cell culture, and found that ATA could drastically inhibit SARS-CoV replication, with viral production being more than 1000 fold than that in the untreated control. ATA is also shown to be an effective anti-viral for several other viruses, including West Nile Virus and variola virus.

BACKGROUND OF THE INVENTION Submersible turbines running on compressed gas, typically air, are known. See for example U.S. Pat. No. 272,656 to Cook, U.S. Pat. No. 271,040 to Cook, and U.S. Pat. No. 211,143 to Fogarty. Since turbines of this type run on low pressure gas, they are suitable for extracting energy from sources (e.g. low pressure gas wells) that are not otherwise useful. Moreover, such turbines are especially suited for use in environments where electric and combustion engines would be unsafe (e.g., in the presence of explosive gases). The compressed air used to run the turbine may also be used to provide necessary ventilation. However, prior art devices present several problems. First, the submerged wheel is often provided with a housing that affords substantial clearance between the wheel and the housing at some point along the circumference. This large clearance results in the escape of air in the form of large bubbles that are released as the turbine wheel rotates. Once the air is clear of the turbine wheel, the upward motion of the air pocket is no longer available for driving the wheel. Moreover, as large amounts of water flow in to replace the air, large scale turbulence is created which extends to the surrounding water. In this manner, the efficiency of the turbine is reduced as energy is dissipated in the surrounding water. The uncontrolled influx of replacement water also causes local turbulence which produces vibrations of the turbine. The vibration both represents a loss of useful mechanical energy output and may even render the turbine output unsuitable for certain applications. Second, typical prior art devices use a wheel submerged within a large tank of water. This results in an extremely heavy machine, possibly unsuitable for installation in existing structures. If a specially designed vessel is used, as for example a sunken well, the turbine cannot readily be moved from one location to another. Aside from the large amount of water that is required to initially put one of these turbines into operation, the large tank represents a substantial cost for material and fabrication. A third difficulty with prior art devices relates to loading on the axle bearings due to the weight of the wheel. Excessive loading leads to frictional losses and possible ultimate failure of the bearings themselves. A fourth difficulty with submersible turbines running on compressed air is the considerable fluid shear generated at the tips of the vanes. This represents a further loss as turbulence is set up in the surrounding water. The use of a closely-fitting housing to reduce these losses tends to localize the shear forces and thus puts an added strain on the turbine wheel. A fifth difficulty encountered with the prior art devices involves frictional losses between the downward moving side of the turbine wheel and the water through which it moves. The turbulence set up causes vibrations which reduces the efficiency of the turbine. Prior art turbines typically use curved, back-swept vanes to cut down viscous drag and to provide more effective air entrapment. This type of vane tends to displace water outward which can cause additional turbulence. A sixth difficulty with the prior art submersible turbines relates to the water replacement when the upwardly moving air leaves the vicinity of the rotating turbine wheel. The inflow of water tends to be relatively non-directional and therefore often acts in a direction opposite that in which the wheel is rotating. Again, the result is a reduction in the energy output of the turbine and increased vibration. SUMMARY OF THE INVENTION A turbine wheel is encased in a closely-fitting cylindrical housing filled with water. The outside of the housing is exposed to the air. The center of the wheel is a hollow cylindrical water-tight shell. Thus, the wheel tends to float, thereby reducing downward loading on the axle bearings. The use of a hollow wheel within a closely fitting housing also results in a rather small total volume of water being used. The circumference of the wheel is fitted with a plurality of chambers defined by vanes, each vane having a radial flat portion that tends to propel water tangentially following the direction of rotation of the wheel. Each vane may also have a backswept portion that improves the air trapping capability of the chambers. Gas, typically air, is injected near the bottom of the wheel, just beyond the bottom of the housing. The buoyant force of the air trapped in the chambers imparts a rotational motion to the wheel. An opening near the top of the housing, surrounded by four upwardly extending vertical walls, allows the air to escape and prevents the displaced water from overflowing. The water displaced by the injected gas circulates in the region between the ends of the vanes and the housing. Efficient displacement may be aided by providing an external conduit. Some of the displaced water passes into the conduit which runs along the outside of the housing, the conduit reintroducing the water near the top of the wheel where the trapped gas is released, thereby aiding the efficient replacement of water. In this fashion, an overall circulation of water is obtained. The circulation leads to less relative motion between the wheel and the water with consequentially decreased friction and turbulence. OBJECTS AND ADVANTAGES OF THE INVENTION It is an object of this invention to disclose a submersible turbine wherein losses due to turbulence in the surrounding water are minimized. This is accomplished by providing a closely-fitting housing surrounding the turbine wheel. An advantage of using a closely-fitting housing is that air which is entrapped in the chambers and then escapes as the wheel rotates is prevented from leaving the vicinity of the turbine vanes by the housing. Thus, the upwardly moving vanes are operating in an environment that is mainly air and thus there is less viscous drag on the turbine wheel due to water at the end of the vanes. Another advantage of using a closely-fitting housing is that the air, by being kept in the vicinity of the turbine vane, is able to exert its lift fully, thereby increasing the efficiency of the turbine. A further advantage of using a closely-fitting housing as the sole means of containing the water for the operation of the turbine is a decrease in the amount of water which leads to a lighter, smaller device. It is an object of this invention to disclose a submersible turbine wherein the load on the axle bearings is decreased. This is accomplished by using a hollow wheel which is buoyed up by the water surrounding it. By adjusting the weight of a given size wheel, neutral buoyancy may be achieved. An advantage of using a hollow wheel is that friction losses in the bearings are minimized and premature bearing failure due to excessive loading is reduced. Moreover, less water is necessary for the operation of the turbine. This results in a turbine of reduced weight. It is another object of this invention to disclose a submersible turbine wherein shear forces at the ends of the downwardly moving turbine vanes are reduced. A close housing is employed, and the wheel is provided with vanes which tend to propel the water in a tangential manner following the direction of rotation. In order to ensure that the downardly moving water at the outer portion of the downwardly moving side of the turbine wheel circulates freely with the wheel, an upward circulation on the upwardly moving side of the wheel is maintained in the region between the ends of the vanes and the inner surface of the housing. An external conduit for the displaced water may be provided, the external conduit leading to a point near the top of the wheel so that the displaced water circulates upwards and is reintroduced. An advantage of using a circulating water mode is that friction between the downwardly moving vanes and the water surrounding them is less than in the case of static water, due to the reduced relative motion. Also, the orderly escape of air near the top of the wheel results in less vibration and turbulence due to inrushing water. An advantage of using an external conduit to take up the displaced water is that the displaced water freely flows through the regularly shaped smooth conduit. The external conduit avoids the problems caused by circulating too much water in the region between the upwardly moving vanes and the housing. Thus excessive velocity of the upwardly moving water, and back pressure on the downwardly moving water is reduced. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is an overall perspective view of the turbine, partially cut away to show the wheel. FIG. 2 is a fragmentary perspective view of the chambers located on the circumference of the turbine wheel. FIG. 3 is a cross-sectional view of the turbine, to show the relationship between the air injectors and the bypass conduit. FIG. 4 is an exploded perspective view of the bottom portion of the housing showing the air injectors and the bypass conduit inlets. FIG. 5 is a fragmentary perspective of an alternate embodiment of the turbine wheel, especially adapted for use without a bypass conduit. DESCRIPTION OF THE PREFERRED EMBODIMENT FIG. 1 is a perspective view showing the major components of the turbine including housing 10 and wheel 20. The outer housing 10 of the turbine is made up of circular side plates 11 spaced apart by central cylindrical shell 12. For a cylindrical shell having a 48-inch inner diameter, 52-inch diameter steel side plates approximately 1/8-inch thick are suitable. They are spaced apart by approximately 14 inches. Flanges 14 are welded to the edges of cylindrical shell 12 to allow the shell to be bolted to the side plates. A 2-inch wide, 1/4-inch thick neoprene gasket is interposed between flanges 14 and side plates 11 to provide a water-tight seal. Approximately 50 bolts hold cylindrical shell 12 to each side plate 11. An opening 16 extends across housing 10, the opening having an axial dimension approximately equal to that of the cylindrical shell 12. Opening 16 is surrounded by vertical side walls 18 parallel to side plates 11, and vertical walls 19 perpendicular to side walls 18, that prevent water from overflowing out of opening 16. In the preferred embodiment, opening 16 extends from a point near the top of the housing to a point approximately 18 inches past the top of the housing. (For convenience, displacements along the circumference of housing 10 or wheel 20 will be assumed to be along the direction of rotation of wheel 20 as shown by arrow 15.) A plurality of baffles may be provided within the upwardly extending chamber defined by walls 18 to aid in the separation of air and water. Outer housing 10 is shown partially cut away to reveal a portion of turbine wheel 20 which is coaxial with housing 10. As seen in FIG. 2, turbine wheel 20 comprises wheel discs 21 and a relatively large cylindrical shell 22, the circumferential surface 23 of which is covered with a plurality of chambers 24. Wheel discs 21 and cylindrical shell 22 are welded together to provide a hollow, water-tight wheel core. Wheel 20 is provided with a 3/4-inch shaft 26 that extends axially in both directions. Shaft 26 is supported from side plates 11 by bearings 28. Wheel discs 21 have a diameter as close to the inner diameter of housing 10 as is practical. A typical clearance is 1/8-inch. In the preferred embodiment, wheel discs 21 extend beyond cylindrical surface 23 by a distance corresponding to the depth of the chambers 24 increased by a specified distance 25, typically 1/2-inch, to provide a circulation region. A chamber depth of approximately 6 inches for a wheel core having an outer diameter of 34.75 inches is suitable. In addition to the wheel disc extensions beyond cylindrical surface 23, chambers 24 are defined by vanes 32. Each vane 32 has a radial portion 34 extending part of the distance from cylindrical surface 23 to the outer edge of side plate 21, and a backswept portion 36 extending outward angularly from the end of straight portion 34 to within distance 25 of the circumference of the wheel discs. For the dimensions of the wheel described above, sections 34 and 36 are both 4 inches by approximately 13 inches (the width of the wheel). There are 30 chambers around the circumference of the wheel. The vanes are made of a relatively thin material such as 1/16-inch iron. Referring to FIGS. 1, 3, and 4, air injectors 40 and bypass conduits 50 can be seen. In the preferred embodiment, there are two air injectors located centrally on cylindrical shell 12 at 5 inches and 10 inches circumferentially past the bottom of the housing. Pipe having an inner diameter of 1/2-inch is suitable. A bypass conduit is on each side of the line between air injectors 40. Each bypass conduit 50 extends from its own bypass inlet 52, up along the outside of housing 10, and to a point above and past the top of housing 10, where it passes into side wall 18. FIG. 4 shows more closely the arrangement of injectors 40, bypass conduits 50, and bypass inlets 52. Each bypass inlet 52 covers an opening 54 in cylindrical shell 12. Each opening 54 is generally rectangular in shape with an overall dimension of 3 inches axially by 10 inches circumferentially. Each opening 54 extends from 0.75 inches from the edge of cylindrical shell 12 to 3.75 inches from the edge, and from 4.25 inches to 14.25 inches past the bottom of cylindrical shell 12. The entire assembly is supported in stand 76 which also prevents side plates 11 from bowing out relative to one another due to water pressure within the housing. A drain 78, typically 1/2-inch, is provided for emptying the water out of housing 10. Having thus set forth the overall construction of the turbine, the operation can be described. In use, housing 10 is filled with water while cylindrical core 22 remains filled with air. Thus a relatively small volume of water is actually used, since the major portion of the housing is filled with the wheel core. Due to the buoyant force on wheel 20, the downward loading on bearing 28 is relieved. When air, as for example, from a compressor, is injected through air injectors 40, it begins to fill those chambers 24 that are nearest the injectors. Water displaced by the injected air flows into intake conduit 52 and out bypass conduit 50. The buoyant force acting on the trapped air in chambers 24 causes the wheel to rotate. As the wheel rotates, the air-filled chambers, upon nearing the top of the wheel, are facing upward and release their air. The released air is replaced by the water contained by walls 18 and 19 and the water from bypass conduits 50. Not all the water displaced by the injected air passes through bypass conduits 50. Some of it circulates upward in the region bounded radially by the ends of vanes 32 and the inner surface of cylindrical shell 12, and axially by the radially outermost portions of wheel discs 21. This is the volume provided by clearance dimension 25. Due to the small clearance between wheel 20 and housing 10, the chambers 24 on the upwardly moving side of the wheel tend to remain filled with air until they reach the top of the wheel. Since they are moving in a region substantially free of water, they are essentially free of viscous drag. The small clearance tends to have a similar effect on the downwardly moving sides of the wheel, where the wheel is operating in a region substantially free of air. The vanes, each with its flat radial portion, tend to propel the water forward in a tangential manner following the rotation of the wheel. Thus, an overall circulation of the water is set up. On the downwardly moving side the circulation is maintained by having the chambers filled with water being propelled forward by the wheel. On the upwardly moving side of the wheel the water has been displaced into bypass conduit 50 and into the peripheral region of the housing where it flows smoothly. The amount of water that circulates in the bypass conduit relative to the amount that circulates in the housing itself depends in part on the chamber depth relative to clearance dimension 25. FIG. 5 is an embodiment of the turbine wheel suitable for use in a turbine where the housing is not provided with a by-pass conduit. Rather than having a single chamber extending axially across the wheel, a plurality of smaller abutting chambers is provided. A given axial row of chambers, including chambers 80a, 80b and 80c, is defined by radial vanes 82 and 84, a plurality of annular segments 86, and outer core surface 23. The entire wheel comprises a plurality of abutting axial rows of this sort. Vanes 82 and 84, and annular segments 86 are thin material such as 1/16-inch galvanized steel. The chambers are typically 2 inches deep, and 2 inches along the circumference of wheel 20. Since these dimensions are small compared to the wheel diameter, the chambers are approximately rectangular and it is convenient to refer to them as such. The axial dimension of the chamber is approximately 2 inches, in which case the chambers are cubical. For the embodiment of the turbine making use of a wheel of this sort, an air injection manifold must be provided with an injector corresponding to each chamber in a given axial row. Thus for a wheel having 7 chambers axially across the wheel, a 7 injector manifold is appropriate. Since no bypass conduit is provided, the turbine of this embodiment relies on setting up the circulation of water in the outermost peripheral region of housing 10, that is the region between the upwardly moving wheel and the inner surface of cylindrical shell 12. The water in this region would tend to be moving faster than the wheel in order to maintain the circulation. The use of multiple chambers in a given axial row rather than a single chamber extending the axial width of the wheel was motivated by a desire to eliminate the sideways turbulence of water which would displace air prematurely up along the circular sides 21 of the wheel. Typical performance for a turbine having the dimensions and construction described herein is an output shaft velocity of 25 rpm with air input from a 120 volt, 1.25 amp compressor having an output of 30 cfm at 3 psi. The output shaft was capable of lifting 200 pounds a distance of 12 inches in 9 seconds. A sprocket and chain arrangement for power transmission has been found preferable to a belt and pulley system since the latter subjects the output shaft to severe sideways forces necessary to provide adequate friction for the belt. The effect of such sideways loading was a decrease of the output rpm by a factor of 2.

A turbine wheel having a central hollow cylindrical water tight shell is encased in a closely fitting cylindrical housing filled with water. The circumference of the wheel carries a plurality of chambers defined by vanes, each vane having a radial flat portion for propelling water tangentially in the direction of rotation of the wheel. Gas is injected near the bottom of the wheel. The buoyant force of trapped gas in the chambers imparts rotation to the wheel, and water displaced by the injected gas circulates in the region between the ends of the vanes and the housing. An external conduit may facilitate efficient replacement by receiving some of the displaced water near the bottom and reintroducing it near the top of the housing.

CROSS REFERENCE TO RELATED APPLICATIONS The present application is a continuation application of U.S. application Ser. No. 15/140,909 filed on Apr. 28, 2016, which is a continuation application of U.S. application Ser. No. 13/697,031 filed on Apr. 29, 2013, now U.S. Pat. No. 9,329,271, which is a National Stage Application of PCT Application No. PCT/EP2011/001662, filed on Apr. 1, 2011, which claims the benefit of U.S. Provisional Patent Application No. 61/362,810, filed on Jul. 9, 2010, and of German Patent Application No. DE 10 2010 020 925.2, filed on May 10, 2010, and which are hereby incorporated by reference. BACKGROUND OF THE INVENTION The invention relates to a method for optically scanning and measuring an environment. Through use of a known method of this kind, a three-dimensional scan is recorded which is then displayed two-dimensionally. Provided that density and extension of the measurement points are smaller than the pixels of the display, a relatively better visual impression is generated if a gap-filling takes place between the measurement points, i.e., if surfaces are generated from the single measurement points. All measurement points can thus be projected onto one plane and be assigned to single pixels. The intermediate pixels of the plane are then filled, for example, by interpolation. SUMMARY OF THE INVENTION According to an embodiment of the present invention, a method, system and computer program product are provided for displaying a plurality of measurement points in three-dimensional space on a two-dimensional plane of a display screen. The method includes projecting the plurality of measurement points onto the two-dimensional plane, the display screen has a plurality of pixels. Each of the measurement points of the plurality of measurement points is assigned to one of the pixels in the plurality of pixels. A depth value is assigned to each of the plurality of pixels that are assigned one of the measurement points of the plurality of measurement points. A first pixel is selected, the first pixel having a first measurement point of the plurality of measurement points assigned to the first pixel, the first pixel having a first depth value assigned to the first pixel. A first side of the first pixel is searched for a second pixel having a second measurement point of the plurality of measurement points assigned to the second pixel, the second pixel having a second depth value assigned to the second pixel. A second side of the first pixel is searched for a third pixel having a third measurement point of the plurality of measurement points assigned to the third pixel, the second side being opposite the first side, the third pixel having a third depth value assigned to the third pixel. It is determined whether the second measurement point and the third measurement point are on a same object plane based at least in part on the second depth value and the third depth value. The first depth value assigned to the first pixel is changed based on determining the second measurement point and the third measurement point are on the same object plane. BRIEF DESCRIPTION OF THE DRAWINGS Embodiments of the present invention are explained in more detail below on the basis of exemplary embodiments illustrated in the drawings, in which FIG. 1 is a schematic illustration of the assignment and filling of the pixels with a view onto the plane, wherein the adjacent pixels are on the same surface; FIG. 2 is a schematic illustration of the assignment and filling of the pixels, according to FIG. 1 , with a view onto the plane; FIG. 3 is a schematic illustration of the assignment and filling of the pixels with a view onto the plane, wherein the adjacent pixels are located on different surfaces; FIG. 4 is a schematic illustration of the assignment and filling of the pixels, according to FIG. 3 , with a view onto the plane; FIG. 5 is a schematic illustration of a laser scanner in the environment including the display device, and FIG. 6 is a partial sectional illustration of the laser scanner. DETAILED DESCRIPTION OF THE INVENTION Referring to the Figures, a laser scanner 10 is provided as a device for optically scanning and measuring the environment of the laser scanner 10 . The laser scanner 10 has a measuring head 12 and a base 14 . The measuring head 12 is mounted on the base 14 as a unit that can be rotated about a vertical axis. The measuring head 12 has a rotary mirror 16 that can be rotated about a horizontal axis. The point of intersection between the two axes of rotation is designated as the center C 10 of the laser scanner 10 . The measuring head 12 also has a light emitter 17 for emitting an emission light beam 18 . The emission light beam 18 may be a laser beam in the range of wave length of approximately 300 to 1600 nm, for example, 790 nm, 905 nm or less than 400 nm, but other electro-magnetic waves having, for example, a greater wave length can be used. The emission light beam 18 is amplitude-modulated with, for example, a sinusoidal or rectangular-waveform modulation signal. The emission light beam 18 is passed from the light emitter 17 onto the rotary mirror 16 where it is deflected and then emitted into the environment. A reception light beam 20 , which is reflected by or otherwise scattered from an object O, is captured again by the rotary mirror 16 , deflected and passed onto a light receiver 21 . The direction of the emission light beam 18 and of the reception light beam 20 results from the angular positions of the rotary mirror 16 and the measuring head 12 , which depend on the positions of their respective rotary drives which are, in turn, detected by a respective encoder. A control and evaluation device 22 has a data link connection to the light emitter 17 and to the light receiver 21 in the measuring head 12 , parts thereof being arranged also outside the measuring head 12 , for example as a computer connected to the base 14 . The control and evaluation device 22 determines, for a multiplicity of measurement points X, the distance d of the laser scanner 10 from the illuminated point on the object O, and from the propagation times of emission light beam 18 and reception light beam 20 . For this purpose, the phase shift between the two light beams 18 and 20 can be determined and evaluated. Through use of the relatively rapid rotation of the mirror 16 , scanning takes place along a circular line. Also, through use of the relatively slow rotation of the measuring head 12 relative to the base 14 , the entire space is gradually scanned with the circular lines. The totality of the measurement points X of such a measurement shall be designated as a scan. The center C 10 of the laser scanner 10 defines for such a scan the origin of the local stationary reference system. The base 14 is stationary in this local stationary reference system. In addition to the distance d to the center C 10 of the laser scanner 10 , each measurement point X comprises a brightness value which may also be determined by the control and evaluation device 22 . The brightness is a gray-tone value which is determined, for example, by integration of the bandpass-filtered and amplified signal of the light receiver 21 over a measuring period which is assigned to the measurement point X. Through use of a color camera, it is optionally possible to generate pictures, by which colors (R, G, B) can be assigned as a value to the measurement points X in addition to the brightness or comprising the brightness. A display device 30 is connected to the control and evaluation device 22 . The display device 30 can be integrated into the laser scanner 10 , for example into the measuring head 12 or into the base 14 , or it can be an external unit, for example part of a computer which is connected to the base 14 . The display device 30 has a graphic card 32 and a screen 34 which can be arranged separately from one another or as a structural unit. The control and evaluation device 22 provides 3D data of the scan. Referring also to FIGS. 1-4 as well as FIGS. 5 and 6 , the graphic card 32 converts the 3-D data into 2-D data (e.g., rendering data), which are displayed on the screen 34 . The 3-D data are the measurement points X, wherein several scans from different positions of the laser scanner 10 can be assembled into one scene. For representing the 2-D data, there are pixels P, i.e., adjacent, small polygonal surfaces (e.g. squares or hexagons), which are arranged in a two-dimensional plane E which corresponds to the screen 34 . The starting point is the projection of the measurement points X onto the plane E with a viewer (e.g., eye, camera), at a certain viewpoint V. The projection appears to be in perspective (i.e., the viewpoint V is in the finite) or orthographical (i.e., the viewpoint V in is the infinite). The projected measurement points X are assigned to single pixels P. A Z-buffer serves for representing the 2-D data, i.e., a two-dimensional auxiliary field for the pixels P. In this Z-buffer, a field element (e.g., depth z) is assigned to each pixel P. The depth z of each projected measurement point X corresponds to the distance of the measurement point X to the plane E with respect to the viewpoint V. The field of the pixels P and the Z-buffer may be treated in the same way as the images. The viewpoint V may be arbitrary per se and is usually changed several times when regarding the scan and/or the scene. Since the measurement points X are punctiform, with gaps in between, and the pixels P usually, in the case of nearby objects O, have a higher density in the plane E than the projections of the measurement points X, a gap-filling is carried out to fill as many pixels P as possible for an improved representation. The graphic card 32 carries this out in parallel using the 3-D data and the indication of the viewpoint V and of the plane E. Initially only those pixels P are filled to which the projection of a measurement point X is assigned, i.e., which exactly cover one measurement point X. These pixels P are filled with the value of the assigned measurement point X, i.e., brightness and, where applicable, color. All other pixels P, which do not exactly correspond with a projection of a measurement point X, i.e., which are “in between” are empty at first, for example are set to zero. Each of the depths z, i.e., the field elements of the Z-buffer, which are assigned to the initially filled pixels P, is set to that depth z 0 , z 1 , z 2 , which corresponds to the distance of the assigned measurement point X to the plane E. All other field elements of the Z-buffer (e.g., depths z) are set to an extreme value, for example, to infinite. If, when the projection of the measurement points X is made, it turns out that two measurement points X are available for one pixel P, the measurement point having the smaller depth z is selected and the other one is rejected, so that covered surfaces and covered edges are not visible. According to embodiments of the present invention, gap-filling takes place in dependence on the depth z 0 , z 1 , z 2 , i.e., on the distance to the plane E. The graphic card 32 selects all pixels P in parallel with respect to time. By way of example, one selected pixel P 0 is regarded now. The assigned depth z, i.e., field element of the Z-buffer, contains the depth z 0 . For each selected pixel P 0 the adjacent pixels P 1 , P 2 , are searched consecutively, i.e., to the left and to the right and above and below. If the adjacent pixel P 1 is not yet filled or if its depth z is bigger than the depth z 0 of the selected pixel P 0 , it is skipped and the second next pixel P is taken as adjacent pixel P 1 , if necessary iteratively. If an adjacent pixel P 1 , the depth z 1 of which is smaller than the depth z 0 of the selected pixel P 0 , is found in one of the directions, a change to the next direction takes place, and it is looked for the adjacent pixel P 2 (e.g., the depth z 2 of which is smaller than the depth z 0 of the selected pixel P 0 ). It is possible to define a maximum number of skipped pixels, i.e., if the adjacent pixel P 1 or P 2 is not yet found after skipping the maximum number of skipped pixels, the search for P 1 or P 2 is aborted. If the adjacent pixels P 1 and P 2 to the selected pixel P 0 have been found in opposite directions, with the depths z 1 and z 2 of the adjacent pixels P 1 and P 2 being smaller than the depth z 0 , it is checked whether P 1 and P 2 are on the same plane, i.e., whether the amount of the difference of z 2 and z 1 is below a threshold value for the depth z crit , i.e., | z 2 −z 1 |<z crit applies. In such a case, the selected pixel P 0 is filled with the value which is interpolated between P 1 and P 2 , i.e., brightness and, if applicable color. The assigned field element of the Z-buffer is likewise set to the interpolated depth between z 1 and z 2 . Interpolation depends on the distance of the selected pixel P 0 from P 1 and P 2 in plane E. If the difference of the depths is too big, i.e., the condition | z 2 −z 1 |>z crit applies, it is assumed that P 1 and P 2 are located on different planes. The selected pixel P 0 is then filled with the value, i.e., brightnesses and, if applicable color, of, for example, the more remote pixel P 1 or P 2 , and the assigned field element of the Z-buffer with the bigger depth z 1 or z 2 . Alternatively, the value and the depth of pixel P 1 or P 2 having the smaller depth z 1 or z 2 is transferred. In the case of more than two adjacent pixels P 1 , P 2 , the average value of the majority, i.e., of the adjacent pixels P 1 , P 2 , which are located on the same surface, can be transferred. Selected pixels P 0 , which are already filled with values of the measurement points, are overwritten by the interpolation of the values of the adjacent pixels P 1 and P 2 . Alternatively, a selected pixel P 0 , which is already filled, remains unvaried. If pixels P have been skipped when finding the pixels P 1 and P 2 , because they were not filled or because their depth z was too big, their adjacent pixels P 1 , P 2 are the same as with the selected pixel P 0 , so that the skipped pixels P and the assigned field elements of the Z-buffer, within the framework of the selection taking place in parallel, are likewise filled either with a value which is interpolated between the pixels P 1 and P 2 and/or the depths z 1 and z 2 (depending on the distance of the selected pixel P 0 from P 1 and P 2 in plane E) or with the value and/or the depth z 1 or z 2 of the more remote one among pixels P 1 or P 2 (or the average value of the majority). Due to the selection taking place in parallel, filling with the value and/or the depth z 1 or z 2 of the more remote among the pixels P 1 or P 2 on account of a difference of depths which is too big, leads to the closer-by pixel P 1 or P 2 forming an edge. Even if no adjacent pixel P 1 or P 2 is found, the depth z 1 or z 2 of which is smaller than the depth z 0 of the selected pixel P 0 , since the selected pixel P 0 is at the side of the screen 34 , an edge is generated, since these selected pixels P 0 at the edge are not filled then. Gap-filling may take place once again to fill further pixels, i.e., to improve the representation once again. Gap-filling may take place in the control and evaluation device 22 or by software running on an external computer. Due to the savings in time by a parallel selection, the hardware-based gap-filling on the graphic card 32 may be used together with the programming interface of the latter.

A method, system and computer program product are provided for displaying three-dimensional measurement points on a two-dimensional plane of a display screen having a plurality of pixels. The method includes projecting the measurement points onto the plane. Each of the measurement points is assigned to one of the pixels. A depth value is assigned to each of the pixels. A first pixel is selected having a first measurement point and a first depth value. A first side is searched for a second pixel having a second measurement point and a second depth value. A second side is searched for a third pixel having a third measurement point and a third depth value. It is determined whether the second and third measurement points are on a same plane. The first depth value of the first pixel is changed when the second and third measurement points are on the same plane.

CROSS REFERENCE TO RELATED APPLICATIONS This Application claims priority of Taiwan Patent Application No. 100108680, filed on Mar. 15, 2011, the entirety of which is incorporated by reference herein. BACKGROUND OF THE INVENTION 1 . Field of the Invention The present invention relates to control methods and devices of a simple node transportation system, and in particular relates to control methods and devices of a simple node transportation system for human-machine interface. 2 . Description of the Related Art The most common example of a simple node transportation system is a common building elevator system. In general, the transportation system includes a plurality of nodes and a vehicle. The vehicle may stop by the node to facilitate loading or unloading of people or goods. A plurality of nodes of the transportation system is usually located on a route. Except for two terminal nodes of the ends of the route, any node between the terminal nodes has two adjacent nodes. Depending on the system requirements, the vehicle may stop by the nodes which have transportation requests. Each system may comprise many transportation routes and vehicles corresponding to transportation routes, those may be controlled by a control center. Take the vertical moving elevator transportation system as an example; each floor with an entrance door of the elevator is regarded as a node. The elevator transportation system may comprise one elevator shaft and an elevator, or many elevator shafts and many elevators. More than one elevator route may share a set of the nodes that the elevators stop by. For example, two elevators both stop at the first floor; one of the two elevators stops at odd floors and the top floor, and the other stops at even floors and the top floor. In another example, the floors at which the two elevators stop at are the same. There are two control types of a simple node transportation system. The first type may be called as intelligent transportation system, which installs a complex control interface at each node. The user can input the target node which he wants to go, and the control center of the system dispatches a vehicle to stop by the node where the user fed input. After the vehicle carries users and/or goods, the control center of the system sends a signal to the vehicle for going to the target node. In transit, the vehicle may stop by other nodes due to other requests, but users and/or goods would only leave the vehicle at the target node. Except for an emergency interface in the vehicle, the vehicle may not be equipped with any control interface. The user only needs to input a command once at the node. Besides, the user does not need to care about the relative direction of the target node. This is why the system called intelligent transportation system. The second type is more traditional, the system installs a simpler control interface in each node, and the user has to determine by himself the direction of the node to which he wants to go and inputs the direction in the control interface. The control center of the system dispatches a vehicle to stop by the node which the user has input. The user has to determine whether it goes the desired direction when vehicular door opens. After the user enters the vehicle, the user has to input the target node to which he wants to go by a complicated control interface in the vehicle. This type of transportation system requires two-stage inputs, wherein the user inputs the direction of the route at the node in the first stage, and inputs the target node in the vehicle in the second stage. In practice, due to the number of nodes on the same route being usually more than the number of the vehicles, only a simple interface is installed in each node with one complicated interface installed in the vehicle, the second type is more economical than the first type. Therefore, the installation number of second type transportation system is greater than the first type actually. Since the development of consumer electronic systems explores in recent years, the electronic systems have made significant progress, and prices have fallen very rapidly. Therefore, there is a need for integrating several features into the aforementioned human-machine interface of a simple node transportation system through electronic systems, such as advertisement, communication, security, monitoring, warnings, and so on. BRIEF SUMMARY OF THE INVENTION In an embodiment, the invention discloses a simple node transportation system, comprising a vehicle traveling on a route, and a traditional control module configured to control the vehicle, and any combination of a node controller and a vehicle controller, wherein the route comprises a plurality of nodes which the vehicle may stop by. The node controller is installed in one of the plurality of nodes, and the vehicle controller is installed in the vehicle. The node controller further comprises a traditional touch module, an input module, a node control module and an output module. The traditional touch module is configured to connect to the traditional control module, and send a control instruction of a user to the traditional control module. The input module is configured to photograph an image in a target area and at least a gesture of the user. The node control module is configured to recognize the gesture, transfer the control instruction corresponding to the gesture, and output to the traditional control module. The output module is configured to display the image in the target area and the control instruction corresponding to the image. The vehicle further comprises: the traditional touch module, the input module, a vehicle control module and the output module. The traditional touch module is configured to connect to the traditional control module, and send the control instruction of the user to the traditional control module. The input module is configured to photograph an interior image in the vehicle and at least a gesture of the user. The vehicle control module is configured to recognize the gesture, transfer the control instruction corresponding to the gesture, and output to the traditional control module. The output module is configured to display the interior image and the control instruction corresponding to the interior image. In another embodiment, the invention discloses an intelligent control module in a simple node transportation system. The simple node transportation system comprises a vehicle traveling on a route, and a traditional control module configured to control the vehicle, wherein the route comprises a plurality of nodes, which the vehicle may stop by. The simple node transportation system comprises at least one of the following components a node controller installed in at least one of the plurality of nodes and a vehicle controller installed in the vehicle. The intelligence control module comprises a network module configured to connect to any combination of a node control module of at least one node controller and a vehicle control module of the vehicle controller, wherein the node control module or the vehicle control module gets an image in the target area photographed by an input module of another node controller or an interior image photographed by the input module of the vehicle controller through the network module, and sends a signal to an output module of the node control module or the vehicle control module to display the image in a target area or the interior image. In another embodiment, the invention discloses a node controller in a simple node transportation system. The simple node transportation system comprises a vehicle traveling on a route, and a traditional control module configured to control the vehicle, wherein the route comprises a plurality of nodes, which the vehicle may stop by. The node controller is installed in at least one of the plurality of nodes, and the node controller comprises: an input module, a node control module and an output module. The input module is configured to photograph an image in a target area and at least a gesture of the user. The node control module is configured to recognize the gesture, transfer the corresponding control instruction and output the control instruction to the traditional control module. The output module is configured to display the image in the target area and the corresponding control instruction. In another embodiment, the invention discloses a vehicle controller in a simple node transportation system. The simple node transportation system comprises a vehicle traveling on a route, and a traditional control module configured to control the vehicle, wherein the route comprises a plurality of nodes, which the vehicle may stop by. The vehicle controller is installed in the vehicle, and the vehicle controller comprises: an input module, a vehicle control module and an output module. The input module is configured to photograph an interior image in the vehicle and at least a gesture of the user. The vehicle control module is configured to recognize the gesture, transfer and output the corresponding control instruction to the traditional control module. The output module is configured to display the interior image in the vehicle and the corresponding control instruction. In another embodiment, the invention discloses a control method of a simple node transportation system. The simple node transportation system comprises a vehicle traveling on a route, and a traditional control module configured to control the vehicle, wherein the route comprises a plurality of nodes, which the vehicle may stop by, and the vehicle controller is installed in the vehicle, and the control method comprises: detecting that at least one user has entered a target area of the node; detecting the gesture of the user in the target area; recognizing the gesture and transferring the corresponding control instruction, and outputting the corresponding control instruction to the traditional control module; and displaying the corresponding control instruction. In another embodiment, a control method of a simple node transportation system is provided. The simple node transportation system comprises a vehicle traveling on a route, and a traditional control module configured to control the vehicle, wherein the route comprises a plurality of nodes, which the vehicle may stop by. The control method comprises: detecting that at least one user has entered a target area of the node; detecting the gesture of the user in the target area; recognizing the gesture and transferring the corresponding control instruction, and outputting the corresponding control instruction to the traditional control module; and displaying the corresponding control instruction. BRIEF DESCRIPTION OF THE DRAWINGS The present invention can be more fully understood by reading the subsequent detailed description and examples with references made to the accompanying drawings, wherein: FIG. 1 is a schematic diagram illustrating an embodiment of a simple node transportation system of the disclosure; and FIG. 2 is a schematic diagram illustrating an embodiment of a node controller of the disclosure; and FIG. 3 is a flowchart of an embodiment of a user control method of the disclosure; and FIG. 4 is a schematic diagram illustrating an embodiment of a vehicle controller of the disclosure; and FIG. 5 is a flowchart of an embodiment of a user control method of the disclosure; and FIG. 6 is a schematic diagram illustrating an embodiment of each component connected in a simple node transportation system of the disclosure; and FIG. 7 is a schematic diagram illustrating an embodiment of each component connected in a simple node transportation system of the disclosure; and FIG. 8 is a schematic diagram illustrating an embodiment of a node controller of the disclosure; and FIG. 9 is a schematic diagram illustrating an embodiment of a vehicle controller of the disclosure; and FIG. 10 is a schematic diagram illustrating an embodiment of each component connected in a simple node transportation system of the disclosure; and FIG. 11 is a schematic diagram illustrating an embodiment of each component connected in a simple node transportation system of the disclosure; and FIG. 12 is a schematic diagram illustrating an embodiment of an intelligent control module of the disclosure. FIGS. 13A and 13B are diagrams illustrating user control of a node controller and a vehicle controller, respectively. FIG. 14 is a diagram illustrating a display in accordance with an embodiment of the invention. FIG. 15 is a diagram illustrating a display in accordance with an embodiment of the invention. FIG. 16 shows a static gesture, such as raising the right hand. FIG. 17 shows a first gesture of raising the right hand or waving the right palm. FIG. 18 shows a first user and a second user inputting two gesture directions simultaneously. FIG. 19 is a schematic diagram illustrating an embodiment of the depth detection module detecting a depth signal of the user in the target area of the disclosure. DETAILED DESCRIPTION OF THE INVENTION Please refer to FIG. 1 . FIG. 1 is a diagram illustrating a simple node transportation system 100 according to an embodiment of the invention. The transportation system 100 comprises two routes 120 and 140 , and two vehicles 122 and 142 . The vehicle 122 travels on the route 120 , and the other vehicle 142 travels on the route 140 . These two routes 120 and 140 may stop by a plurality of nodes 160 , which include two terminal nodes 160 a and 160 b . As shown, the vehicle 142 stops by the terminal node 160 a , and the vehicle 122 stops by the node 160 . In an example, the sets of the nodes 160 which the two routes 120 and 140 stop by are the same. In another embodiment, the sets of the nodes 160 which the two routes 120 and 140 stop by are different, but at least share one common node 160 . Each node 160 is equipped with a node controller 162 as a human-machine interaction interface. The terminal nodes 160 a and 160 b are equipped with the terminal node controllers 162 a and 162 b . The vehicles 122 and 142 are equipped with a respective vehicle controller 124 and 144 as a human-machine interaction interface. The node controller 162 and the vehicle controllers 124 and 144 are connected to a control device 110 . The control device 110 controls the vehicle controllers 124 and 144 according to the instructions received from the human-machine interface; and the control device 110 demands the vehicle controllers 124 and 144 to travel between the nodes 160 on the route and stop by the nodes 160 to load and unload people and goods. The nodes 160 and the vehicles 122 and 142 may be equipped with security doors (not shown), the control device 110 may also control the security doors for opening/closing. Please refer to FIG. 2 . FIG. 2 is a diagram illustrating a node controller 162 according to an embodiment of the invention. The node controller 162 comprises a traditional touch module 210 , an input module 220 , an output module 230 , and a node control module 240 . The traditional touch module 210 comprises a panel and the buttons with the direction indicator lights 212 , and 214 . The operation mode of the traditional touch module 210 is similar to the second type described in the description of the related art. The user first determines which direction of the target node he wants to go toward, and presses the button with the direction indicator light 212 or 214 corresponding to the direction, and then the button with the direction indicator light 212 or 214 lights up. After the vehicle arrives at the target node and opens the security doors, the button with the direction indicator light 212 , or 214 goes off. User keeps pressing the button on the direction indicator light 212 or 214 to make the security doors of the vehicle keep open. The input module 220 and the output module 230 are connected to the node control module 240 . The input module 220 may comprise a mounting assembly 222 to attach the input module 220 to a proper location. The mounting assembly 222 may comprise a control mechanical equipment such that the whole input module 220 could be pitched and/or rotated in one-dimensional or multi-dimensional degrees of freedom. The input module 220 may comprises one or more of sound reception module 224 to receive monaural or multi-channel stereophonic sound. When the received volume is larger than a threshold value, the sound reception module 224 may send a signal to activate the whole or a part of the node controller 162 . If the sound reception module 224 does not receive sound over a certain volume within a certain time period, the whole or a part of the node controller 162 may switch to energy-saving mode that may save more electricity. The input module 220 may comprises one or more of depth detection modules 226 to detect distance of presented object in front of the input module 220 . The depth detection module 226 may be implemented in various manners including photographic lens that has multiple overlapping angles of vision, a laser rangefinder, an ultrasonic distance measurement device and so on. The present invention does not limit implementation choices of the depth detection module 226 , as long as the implementation is capable to identify the distance between the object and the input module 220 . As shown in FIG. 19 , FIG. 19 is a schematic diagram illustrating an embodiment of the depth detection module detecting a depth signal of the user in the target area of the disclosure. The input module 220 may comprise a photographic module 227 and a lighting module 228 , wherein the lighting module 228 may emit the wavelength of visible light and the wavelength of infrared ray to illuminate the target area. The photographic module 227 may photograph the images in the wavelength of visible light and of infrared ray. Because there may be many complicated lighting situation for the target area photographed by the photographic module 227 , the multi-spectral photography may filter out the noise to get clearer images. The photographic module 227 may also have capability for zooming out or zooming in. The input module 220 further may comprise a motion detection module 229 . When object goes into or through the target area, the motion detection module 229 may send a signal to activate all or a part of the whole input module 220 or the node controller 162 to start via the node control module 240 . However, if no object goes into or through the target area within a certain period of time, all or a part of the node controller 162 may switch to the energy-saving mode that may save more electricity. The depth detection module 226 may delimit the interesting distance between the target area and the input module 220 . In an embodiment, if the target area is an open area in front of the input module 220 , there may be many people walking around in the target area. If the input module 220 only uses the photographic module 227 and/or the motion detection module 229 , the target area may be too large. Therefore, the depth detection module 226 may be configured to restrain the depth of the target area to avoid the misjudgment resulting from the object moving behind the target area. The node control module 240 is configured to receive the signal input from each module in the input module 220 , and further may process and output the signal. The part of signal processing may comprise at least three levels. The first level may comprise signal sampling, compression, format conversion, storage, and re-output. For example, the sound reception module 224 outputs the signal to the node control module 240 , and the node control module 240 may perform sampling, compression, format conversion, storage, and re-output of the audio signal. The photographic module 227 outputs the signal to the node control module 240 , and the node control module 240 may perform sampling, compression, format conversion, storage, and re-output of the video signal. The depth detection module 226 outputs the signal to the node control module 240 , and the node control module 240 may perform sampling, compression, format conversion, storage, and re-output of the depth signal. The second level is the node control module 240 performing data fusion or integration between different media or related processing. For example, the node control module 240 laps the video signal over the depth signal or performs the related processing, and outputs a three-dimensional video signal. In addition, the node control module 240 may integrate the video signal with the depth signal or performs the related processing, and outputs a three-dimensional video signal animation. The third level of the signal processing involves the recognition of the media content; especially when the node control module 240 uses the data fusion or integration of two or more media or related processing to recognize people and gesture in the target area. When the signal output from the photographic module 227 and/or the signal output from the depth detection module 226 are integrated into the output signal, the node control module 240 at least can perform face recognition and gesture recognition for people. Face recognition comprises at least position recognizing art and characteristic recognizing art. The characteristic recognizing art used to recognize the gender, the identity or the approximate age of the user through the characteristics of faces is more sophisticated than the position recognizing art used to recognize the positions where faces are. In additional to recognizing so-called hard point such as the center of a palm and/or finger tips, the recognition may include the vertices of the face, elbows, shoulders, neck, hips, knees and other vertices. The lines between corresponding vertices and recognized human parts form a skeleton of a human body. With the change of the time axis, the node control module 240 may recognize the movements of the human body, such as raising hands, waving hands, shaking hands and so on. According to the node control module 240 in this embodiment of the invention, the node control module 240 may perform the three levels of the signal processing. The output module 230 comprises a display 231 and more than one speaker 232 . The display 231 may show separate windows, which comprise a first target area window 234 configured to display the situation of the target area corresponding to the node controller 162 . Next, please refer to the FIG. 3 , which is a flowchart of a user control method 300 according to an embodiment of the invention. The user may control the node controller 162 installed in the node 160 according the method 300 , as shown in FIG. 13A . In step 310 , the node control module 240 may detect that one or many users have entered the target area where the node controller 162 is located by the sound reception module 224 , the motion detection module 229 and/or the photographic module 227 . Next, in an optional step 320 , the node control module 240 may adjust the lighting module 228 for lighting the target area to make the photographic module 227 photograph the illuminated target area clearly. The node control module 240 may also adjust the direction of the mounting assembly 222 to make the photographic module 227 focus on the user in the target area clearly. In step 330 , the node control module 240 sends the image photographed by the photographic module 227 to the first target area window 234 of the display 231 . Furthermore, the node control module 240 may also remind the user that he has entered the target area where the node controller 162 is located by the speaker 232 . In an embodiment, the node control module 240 may recognize the characteristics of the user and mark the characteristics according to the three levels of the signal processing. The manners of marking the characteristics may comprise but not be limited to the following manners: framing the human face; displaying the user's ID or name if the node control module 240 has recognized the identity of the user; and/or verbally greeting the user by user's ID or name. For example, the node control module 240 produces a voice of “someone, hello, may I ask you to go upstairs or downstairs?” through the speaker 232 ; marks the control hard points like the palm of hands /fists /fingers and so on; and marks the vertices of each joint of the human body and the lines between corresponding vertices. Next, in step 340 , the node control module 240 detects the first gesture of the user to start the control process of the transportation system. The first gesture mentioned in the invention may comprise a static gesture, such as raising the right hand in FIG. 16 , and may also comprise a dynamic action, such as waving the right palm, holding or opening a fist, and may also comprise static gestures and dynamic actions, such as raising the right hand and doing the pose of holding or opening the fist. In step 330 , the node control module 240 may also prompt the user to do the first gesture by the sound or the image, and makes the user that uses the node controller 162 at the first time or are not familiar with the operation of the node controller 162 operate smoothly. In one embodiment, the first gesture may comprise plural kinds of gestures and/or actions, such as raising the right hand or waving the right palm in FIG. 17 . As long as the user does one of the gestures, the node control module 240 regards the gesture as the first gesture. Similarly, after detecting that the user does the first gesture, the node control module 240 may frame the user specially, make a sound to confirm that the user has inputted the first gesture, and guide the user to do the second gesture. In an embodiment, the second gesture may comprise the plural kinds of gestures and/or actions. For example, the gesture of turning the palm up and the gesture of turning the palm down would correspond to the two directions of the movement of the vehicle respectively. In another embodiment, the second gesture may comprise the positions in which the control hard points of the user are showed in the display 231 . For example, the user moves the control hard points to the direction control area 233 in the display 231 within a period of time or does the gesture of clenching the fist/finger splay. In step 350 , the node control module 240 detects the second gesture of the user to arrange for the vehicle. The effect is like pressing the button with the direction indicator light 212 , or 214 of the traditional touch module 210 , wherein the operation mode is the second type described in accordance with the prior art. Finally, in step 360 , the node control module 240 may turn on the light corresponding to the direction in the direction control area 233 of the display 231 , and may also make a sound to confirm the direction, and may turn on the button with the direction indicator light 212 , or 214 corresponding to the direction in the traditional touch module 210 , as shown in FIG. 14 . In a similar example, if the user wants to cancel the previous instructions, the node control module 240 may detect the third gesture of the user, and then turn off the light corresponding to the direction light in the direction control area 233 of the display 231 , and also make a sound to confirm the direction, and turn off the button with the direction indicator light 212 , or 214 corresponding to the direction in the traditional touch module 210 . In another embodiment, many users can operate the node controller 162 simultaneously. For example, the first user and the second user may operate the operation of inputting two directions simultaneously, as long as the target area may accommodate many users, and the node controller 162 may analyze the gestures and actions. In other words, the input method 300 used by two or more users may be used in different steps. For example, when the input method 300 used by the first user stays in step 330 , the input method 300 used by the second user may proceed as in the step 350 . Back to FIG. 2 , in many windows of the display 231 , in addition to the direction control area 233 and the first area window 234 introduced before, the display 231 may comprise a second target area window 236 , an advertisement area window 238 and an emergency notification area window 239 . The second target area window 236 is configured to display the node which the vehicle stops by, and the situation of the target area of the node. In FIG. 1 , for example, assume that the movement direction of the vehicle 122 is downward, and the movement direction of the vehicle 142 is upward. In any node 160 or 160 b , after the user inputs the instruction that the direction is upward, the second target area windows 236 of these node controllers 162 display information that the node that the vehicle 142 stopped by is 160 a and the audio and video recorded by the input module 220 of the node controller 162 a . When the vehicle 142 does not appear, the user in each node can monitor the current situations of the vehicle 142 through the second target area window 236 . The advertisement area window 238 may broadcast the wireless television programs, the programs stored in advance, a temporary scrolling text marquee advertisement and so on. Furthermore, the advertisement area window 238 may interact with the user by playing the simple gesture game. For example, stretching exercises, throwing or catching a ball, dancing and so on. As long as the user does not use the first and second gestures of the transportation system in the game, the node controller 162 may even allow the user to play the game and operate the control method of the transportation systems at the same time. Finally, the emergency notification area window 239 is configured to allow the user to start the emergency notification area window 239 through an emergency gesture when the user encounters an emergency. The emergency gesture may be a “full time” gesture. It means that no matter when it is, as long as the node control module 240 detects that any person in the target area does this emergency gesture, then the node control module 240 enters the situation of the emergency notification. In another embodiment, as long as the node control module 240 detects that the hard points of the user have moved into the emergency notification area 239 and the emergency gesture is formed by the hard points, the node control module 240 would enter the situation of the emergency notification. After the node control module 240 enters the situation of the emergency notification, the user can talk to the handler who deals with the emergency through the sound reception module 224 of the input module 220 , the photographic module 227 and the output module 230 . In the situation of the emergency notification, the node control module 240 records and stores the audio, video, and even the depth of the signal for retrieving the records in the aftermath. It is noted that, although FIG. 3 shows that the control method 300 includes step 340 and step 350 , wherein step 340 and step 350 detect the first gesture and the second gesture of the user respectively, the reason of detecting the first and the second gesture is to reduce the probability of the misjudgment. In another embodiment of the invention, after detecting directly the second gesture of the user, the node control module 240 may arrange for the vehicle immediately. Please refer to FIG. 4 , which shows a diagram illustrating a vehicle controller 124 according to an embodiment of the invention. The vehicle controller 124 or 144 comprises four modules, which are a traditional touch module 410 , an input module 220 , an output module 230 and a vehicle control module 440 , respectively. The input module 220 and the output module 230 are connected to the vehicle control module 440 respectively. The traditional touch module 410 comprises a panel and a plurality of buttons with the direction indicator lights. In this example, a plurality of nodes represent the first floor to the sixth floor respectively, and therefore 1F˜6F represent the first floor to the sixth floor. The operation mode of the traditional touch module 410 is similar to the second type described in the description of the related art, the user first determines which direction of the target node he wants to go toward, and presses the button with the direction indicator light corresponding to the direction, and then the button with the direction indicator light lights up. The input module 220 of the vehicle controller 124 and the input module 220 of the node controller 162 are the same basically, so the input module 220 is not mentioned here. The output module 230 of the vehicle controller 124 and the output module 230 of the node controller 162 are the same basically. The different part is that the direction control area 233 of the display 231 is changed to a node indicating area 432 . The node indicating area 432 displays the node corresponding to the traditional touch module 410 . In the above example, the node indicating area 432 shows that six nodes represent the first to the sixth floor respectively. Please refer to FIG. 5 , which shows a diagram illustrating a user control method 500 according to an embodiment of the invention. The user may control the vehicle controller 124 installed in the vehicle 122 according the method 500 , as shown in FIG. 13B . The control method 500 is quite similar to the control method 300 , and for most of the control method 500 may be referenced to the steps of the control method 300 . In step 510 , the vehicle control module 440 may detect that one or many users have entered the vehicle by the sound reception module 224 , the motion detection module 229 and/or the photographic module 227 . Next, in an optional step 520 , the vehicle control module 440 may adjust the lighting module 228 for lighting the target area to make the photographic module 227 photograph the illuminated target area clearly. The vehicle control module 440 may also adjust the direction of the mounting assembly 222 to make the photographic module 227 focus on the users in the target area clearly. In step 530 , the vehicle control module 440 sends the image photographed by the photographic module 227 to the first target area window 234 of the display 231 . Furthermore, the vehicle control module 440 may also remind the user that he has entered the vehicle through the speaker 232 and let the user determine whether he needs to control the vehicle or not. In an embodiment, the vehicle control module 440 may recognize the characteristics and mark the characteristics according to the three levels of the signal processing. The manners of marking the characteristics may comprise but are not be limited to the following several manners: framing the human face; displaying the user's ID or name if the node control module 240 has recognized the identity of the user, and/or verbally greeting to the user's ID or name. For example, the node control module 240 makes a sound of “someone, hello, which floor are you going to?”; marks the control hard points like the palm of hands/fist /fingers and so on; and marks the vertices of each joint of the human body and the lines between corresponding vertices. Next, in step 540 , the vehicle control module 440 detects the third gesture of the user to start the control process of the transportation system. In the step 530 , the vehicle control module 440 may also prompt the user to do the third gesture by the sound or image. Similarly, after detecting that the user does the third gesture, the vehicle control module 440 may frame the user specially, make a sound to confirm that the user has entered the third gesture, and guide the user to do the fourth gesture. In an embodiment, the third gesture and the first gesture may be the same. In an embodiment, the fourth gesture may comprise the plural kinds of gestures and/or actions. For example, the display 231 displays the route and the plurality of nodes. The user turns the palm left and turns the palm right corresponding to two directions of the movement of the vehicle, and the vehicle control module 440 may use the control hard points of the palm to choose the direction in which the user wants to go. In another embodiment, the fourth gesture may comprise the control hard points of the user in the position in the display 231 . For example, the user moves the control hard points to the node indicating area 432 in the display 231 within a period of time or does the gesture of clenching the fists/finger splay, and may use the control hard points of the palm to choose the direction in which the user wants to go. In step 550 , the vehicle control module 440 detects the fourth gesture of the user to arrange for the vehicle. The effect is like pressing the buttons with the direction indicator lights of the traditional touch module 410 , and the operation mode is the second type described in accordance with the prior art. Finally, in step 560 , the vehicle control module 440 may turn on the light of the target node in the node indicating area 432 , may also make a sound to confirm, and may turn on the corresponding button with the direction indicator light in the traditional touch module 410 , as shown in FIG. 15 . In a similar example, if the user wants to cancel the previous instructions, the vehicle control module 440 may detect the third gesture of the user, and then turn off the light corresponding to the target node in the node indicating area 432 of the display 231 , and also make a sound to confirm, and turn off the corresponding buttons with the direction indicator lights in the traditional touch module 410 . It is noted that, although FIG. 5 shows that the control method 500 includes step 540 and step 550 , wherein step 540 and step 550 detect the third gesture and the fourth gesture of the user respectively, the reason of detecting the first and the second gesture is to reduce the probability of the misjudgments. In another embodiment of the invention, after detecting directly the fourth gesture of the user, the vehicle control module 440 may arrange for the vehicle immediately. In another embodiment, many users can operate the vehicle controller 124 simultaneously. For example, the first user and the second user operate the operation of inputting two directions simultaneously in FIG. 18 , as long as the target area may accommodate many users, and the vehicle controller 124 may analyze the gestures and actions. In other words, the input method 500 used by two or more users may be in different steps. For example, when the input method 500 used by the first user stays in step 530 , the input method 500 used by the second user may proceed to the step 550 . In many cases, the user attempts to enter or exit the vehicle while the security doors are closing. In general, although the security doors may be equipped with the security measures to avoid jamming people or goods, taking multiple security measures for the vehicle is still needed to keep safer. According to an embodiment of the invention, the vehicle controller 124 and the node controller 162 may set a prohibited area within a certain range from the security doors. When the security doors is closing, the photographic module 227 and/or the depth detection module 226 of the input module 220 detect that the object is in the prohibited area, and the vehicle controller 124 and the node controller 162 may open the security doors, and may also send a signal to the display 231 and the speaker 232 to issue a warning. Although many simple node transportation systems exist in the world, the control part of the simple node transportation systems still belongs to the traditional type. According to an embodiment of the invention, minimal modification of the original simple node transportation system can be achieved. Please refer to FIG. 6 , which shows a diagram illustrating each component connected in a simple node transportation system 600 according to an embodiment of the invention. The transportation system 600 comprises a control device 110 . The control device 100 further comprises a traditional control module 610 and an intelligent control module 620 . The traditional control module 610 is configured to the traditional touch module 410 of the vehicle controller 124 and the traditional touch module 210 of the node controller 162 . The traditional control module 610 receives the input from the user of two traditional touch modules 210 and 410 , and may control the scheduling and the running of the vehicle. The traditional control module 610 may be configured to connect to a traditional network control center 640 to transmit the running situation of the transportation system 600 to the traditional network control center 640 . In this embodiment, the intelligent control module 620 is configured to the vehicle control module 440 of the vehicle controller 124 and the node control module 240 of the node controller 162 . In an example, the connected-state may present a shape of a star, and the intelligent control module 620 is the center of the star, such that each vehicle control module 440 and each node control module 240 are connected to each other through the intelligent control module 620 . In another example, each component is connected to each other through a bus or the Internet. No matter what the connections, each vehicle control module 440 and each node control module 240 may transmit the signals to each other, and the intelligent control module 620 may be also connected to each vehicle control module 440 and each node control module 240 . The intelligent control module 620 may also be connected to an intelligent network control center 630 to receive the control signal of the intelligent network control center 630 . In this embodiment, the traditional touch module 210 of the node controller 162 and the node control module 240 are connected to each other. After receiving the input from the user, the node control module 240 gives the instruction to the corresponding direction of the traditional touch module 210 through the connecting circuit. After the traditional touch module 210 receives the instruction sent from the node control module 240 , for example, the button of “up stairs” and “down stairs”, the traditional touch module 210 follows the steps to inform the traditional control module 610 , and then the traditional control module 610 plans a schedule for the vehicle. If the user gives the instruction to the button with the direction indicator light 212 , or 214 of the traditional touch module 210 , the node control module 240 also receives a signal indicating what instruction was given by the user through the connecting circuit, and further turns on the light corresponding to the direction of the direction control area 233 in the display 231 . If the user cancels the instruction to the traditional touch module 210 , the node control module 240 also receives a signal indicating what instruction was cancelled by the user through the connecting circuit, and further turns off the light corresponding to the direction of the direction control area 233 in the display 231 . When the vehicle arranged by the traditional control module 610 arrives at the node 160 , the traditional control module 610 turns off the button with the direction indicator light 212 , or 214 of the traditional touch module 210 . When receiving the signal indicating that the button with the direction indicator light 212 , or 214 is turned off through the connecting circuit, the node control module 240 may receive the signal indicating that the vehicle has arrived at the node 160 . Therefore, the node control module 240 may turn off the light corresponding to the direction within the direction control area 233 of the display 231 , and also inform the intelligent control module 620 that the vehicle has arrived at the node 160 . The intelligent control module 620 may inform the node control module 240 of another node 160 , and send a signal to the second target 236 of the display 231 to display the video signal of the input module 210 of the node which the vehicle stops by. The intelligent control module 620 may also inform the vehicle control module 440 , and send a signal to the second target 236 of the display 231 to display the video signal of the input module 210 of the node which the vehicle stops by. The intelligent control module 620 further may inform the intelligent network control center 630 to monitor the signal of the vehicle and the input module 210 of the node which the vehicle stops by. Similarly, in this embodiment, the traditional touch module 410 and the vehicle control module 440 of the vehicle controller 124 are connected to each other. After receiving the input from the user, the vehicle control module 440 gives the instruction to the node corresponding to the traditional touch module 410 by the connecting circuit. After the traditional touch module 410 receives the instruction from the vehicle control module 440 , for example, after pressing the button of “the first floor”, the traditional touch module 410 follows the steps to inform the traditional control module 610 , and then the traditional control module 610 plans a schedule for the vehicle. If the user gives the instruction to the traditional touch module 410 , the vehicle control module 440 also receives a signal indicating what instruction was given by the user through the connecting circuit, and further turns on the light corresponding to the node of the node instruction area 432 in the display 231 . If the user cancels the instruction to the traditional touch module 410 , the vehicle control module 440 also receives a signal indicating what instruction was cancelled by the user through the connecting circuit, and further turns off the light corresponding to the node of the direction control area 432 in the display 231 . When the vehicle arranged by the traditional control module 610 arrives at a certain node 160 , the traditional control module 610 turns off the light of the node of the traditional touch module 410 . When receiving the signal indicating the light of the node is turned off by the connecting circuit, the vehicle control module 440 may receive the signal indicating that the vehicle has arrived at the node 160 . Therefore, the vehicle control module 440 may turn off the light corresponding to the direction of the node instruction area 432 in the display 231 , and also inform the intelligent control module 620 that the vehicle has arrived at the node 160 . Because the simple node transportation system 600 may affect the safety of the passengers, the traditional control module 610 has to be authenticated and testes repeatedly. The advantage showed by an embodiment of FIG. 6 is that the new components do not need to be changed or be connected directly to the traditional control module 610 . The traditional control module 610 does not have to authenticate and test the security function of the core again. However, the con is that the integration level may be lower and the system reaction may be slower. If the simple node transportation system 600 needs a higher degree of integration functions and faster reaction velocity, the simple node transportation system 600 may use the following mode of the connection. In addition, in the embodiment of FIG. 6 , the simple node transportation system 600 does not install the intelligent control module 620 . In an example, as a result of the passengers discharging at some certain nodes 160 being larger, the node controller 162 are merely needed to attach to some nodes 160 for performing the control function by the gestures. In another example, only the vehicle controller 124 may be installed in the simple node transportation system 600 . In other words, the vehicle controller 124 , the node controller 162 , and the intelligent control module 620 may exist alone, or cooperate with each other. Please refer to FIG. 7 , which shows a diagram illustrating each component connected in a simple node transportation system 700 according to an embodiment of the invention. The main difference between the embodiment of FIG. 7 and FIG. 6 is that the traditional control module 610 of the transportation system 700 and the intelligent control module 620 are connected to each other, while the transportation system 600 is connected to the vehicle and the node. In this embodiment, the traditional control module 610 may output and input the position of the vehicle and the signal indicating that the situation that the users issue instructions in each node. Therefore, the intelligent control module 620 may transmit the information transmitted from the traditional control module 610 to each vehicle controller 124 and each node controller 162 , so that the vehicle control module 440 and the node control module 240 corresponding to each vehicle controller 124 and each node controller 162 may be displayed in the node instruction area 432 and the direction control area 233 of the display 231 correctly. The intelligent control module 620 may also transmit the instruction to the traditional control module 610 to arrange for the vehicle according to the input of each vehicle controller 124 and each node controller 162 . It is noted that, in the embodiment of the simple node transportation system 700 , the intelligent control module 620 has to be there, but not every vehicle and node have to be installed the vehicle controller 124 and the node controller 162 . In addition, the touch module 410 of the traditional vehicle controller 124 and the vehicle control module 440 may also be connected to each other. The traditional touch module 210 of the node controller 162 and the node control module 240 may also be connected to each other. The simple node transportation system 700 may also comprises a network control center 710 , which is connected to the traditional control module 610 and the intelligent control module 620 . The network control center 710 may monitor the vehicle and the signal of the input module 210 of the node which the vehicle stops by. The parts described above all improve the controlling mode of the second form described in the prior art, and the following parts modify the intelligent controlling mode of the first form described in the prior art. The control method of the first form is that the user can input the node 160 which the user wants to go to at the node 160 in advance. After entering the vehicle, the user does not input the node 160 which the user wants to go to. Please refer to FIG. 8 , which shows a diagram illustrating a node controller 162 according to an embodiment of the invention. The node controller 162 of the embodiment is actually very similar to the vehicle controller 124 showed in FIG. 4 . The user can choose the node 160 which the user wants to go to through the traditional touch module 410 , and also can choose the node 160 which the user wants to go to through the input module 220 and the output module 230 controlled by a intelligent node control module 840 . After seeing the previous introduction, a person of ordinary skill should be able to understand the operation mode, and therefore the operation mode is not detailed here. In the regulation of teaching in some religion, the users which are a different gender can not take the same vehicle. To avoid the problem of sexual harassment between different genders, in an embodiment of the invention, the intelligent node control module 840 may recognize the gender of the user, and further inform the control device 110 . Therefore, the control device 110 may arrange for a vehicle to stop by the node 160 appropriately. In another embodiment of the invention, if the genders of the group of the users are different, the intelligent node control module 840 may inform the vehicle that is coming or has stopped by the node 160 contains which gender. If the user of the other gender wants to enter the vehicle, the intelligent node control module 840 may issue a warning and/or notify the remote administrator. Please refer to FIG. 9 , which shows a diagram illustrating a vehicle controller 124 according to an embodiment of the invention. In the transportation system of the first form, after entering the vehicle, the user does not have to input the node 160 which he wants to go to again, so the vehicle controller 124 is not equipped with the traditional touch module. The vehicle controller 124 of the embodiment is actually very similar to the vehicle controller 124 in FIG. 4 . The vehicle control module 940 is also very similar to the vehicle control module 440 . The difference is that the vehicle control module 940 does not have to be connected to the traditional touch module 410 . After seeing the previous introduction, a person of ordinary skill should be able to understand the operation mode, and therefore the operation mode is not detailed here. Similarly, if the vehicle has been set for a specific gender, the vehicle control module 940 detects that the user of the other gender has entered the vehicle, and the vehicle control module 940 may issue a warning and/or notify the remote administrator. Please refer to FIG. 10 , which shows a diagram illustrating each component connected in a simple node transportation system 1000 according to an embodiment of the invention. FIG. 10 is similar to FIG. 6 , wherein the intelligent node control module 840 and the traditional touch module 410 are connected to each other, and there is no traditional touch module in the vehicle. Please refer to FIG. 11 , which shows a diagram illustrating each component connected in a simple node transportation system 1100 according to an embodiment of the invention. FIG. 11 is similar to FIG. 7 , wherein the traditional control module 610 of the transportation system 1100 and the intelligent control module 620 are connected to each other. In an embodiment of the invention, the traditional control module 610 and the intelligent control module 620 supply power respectively. In another embodiment, the power of the traditional touch module 210 and 410 and the traditional control module 610 belong to the same system. In addition to the traditional touch module 210 and 410 , the power of each vehicle controller 124 and each node controller 162 and the intelligent control module 620 belong to the same system. Though, these two power systems may be equipped with the uninterruptible power supply device, only one of them has a problem, and the problem does not affect other components of the power system. Please refer to FIG. 12 , which shows a diagram illustrating an intelligent control module 620 according to an embodiment of the invention. The intelligent control module 620 may comprise a network module 1210 configured to be connected to each vehicle controller 124 and each node controller 162 . In an embodiment, the network module 1210 may be connected to an intelligent network control center 630 , the network control center 710 and/or the traditional control module 610 . The intelligent control module 620 may comprise an advertisement module 1220 configured to store various advertisement videos, or receive the signal from other radio broadcast stations for supplying each vehicle controller 124 and each node controller 162 to broadcast in the advertisement area 238 in the display 231 . Because each vehicle controller 124 and each node controller 124 may return the image of the user, in an embodiment, the advertisement module 1220 may provide the individual differentiated advertisements according to the user photographed by each vehicle controller 124 and each node controller 162 . For example, if the user photographed by a certain node controller 162 is a woman, the advertisement module 1220 may send a signal to the node controller 162 to broadcast the advertisements about cosmetics or costume. If the user photographed by a certain node controller 162 is a man, the advertisement module 1220 may send a signal to the node controller 162 to broadcast the advertisements about cameras or computers. The intelligent control module 620 may comprises a record module 1230 configured to store various signals recorded by each vehicle controller 124 and each node controller 162 , wherein the signals include video signals, audio signals and depth signals. The record module 1230 may determine whether the record module 1230 records the signals or not according to the signals transmitted from the motion detection module 229 of the input module 220 . If the motion detection module 229 does not detect any movement, the record module 1230 may not have to record the signal of the vehicle and the node. The intelligent control module 620 may comprises an switch module 1240 configured to exchange the information with the traditional control module 610 . In the embodiment shown in FIG. 7 and FIG. 11 , the intelligent control module 620 may exchange the information with the traditional control module 610 , and the switch module 1240 is configured to provide the interpretation for transforming signals. In addition, the switch module 1240 may also be configured to connect to any vehicle/node and the intelligent network control center 630 /the network control center 710 to let the remote administrator in the intelligent network control center 630 /the network control center 710 communicate with the user in the vehicle/node directly. In conclusion, in addition to the advantages mentioned above, the present invention may also provide at least the following advantages. First, the user does not need to touch any button, and the user can control the simple node transportation system. For example, in health sensitive environments or in hospitals or laboratories needing high infection control, the user can avoid the contact of infection. Furthermore, the simple node transportation system may be inputted with instructions by many users at the same time according to the invention, and does not force the users to squeeze before the conventional touch panel. Especially in a narrow vehicle, the user only lifts a finger, and the user can operate the vehicle. Moreover, the invention may enhance the rate of the attention of the advertisement. The user needs to peer at the display to control the simple node transportation system, and therefore the advertisement in the display may gain higher rate of the attention. If the advertisement is integrated with segment advertisements which may be classified according the passengers, the effect is better than other advertisement machines. Furthermore, the user can see the running situations of the vehicle through the second target area. For example, the current situations of the users pass in and out the vehicle that stops by the node, or the interior situations of the vehicle. The disclosure may provide a better system to know the situations, and avoid users from waiting for the vehicle in the case where they do not know what happened. Finally, the invention may enhance the safety of the system. For example, the setting of the prohibited area may add a layer of insurance on the switches of the security doors, or for example, the system may record the signal at any time and send it to the remote for storing the signal, so that the remote manager can communicate with the user in the node target areas. The above-mentioned examples may enhance the safety of the simple node transportation system.

A simple node transportation system having a vehicle traveling on a route having nodes, a traditional control module controlling the vehicle, and any combination of a node controller and a vehicle controller. The node controller has a traditional touch module connecting to the traditional control module and sending a control instruction to the traditional control module, an input module photographing an image and a gesture, a node control module recognizing the gesture and transferring the corresponding control instruction, and an output module displaying the image and the corresponding control instruction. The vehicle controller has a traditional touch module connecting to the traditional control module and sending a control instruction to the traditional control module, an input module photographing an image and a gesture, a vehicle control module recognizing the gesture and transferring the corresponding control instruction, and an output module displaying the image and the corresponding control instruction.

FIELD OF THE INVENTION [0001] The present invention relates to an electrothermal element, and more particularly to an electrothermal element with better infrared emissivity, which can be used as a heat source of heater and installed in the ceiling or mounted directly on the wall. BACKGROUND OF THE INVENTION [0002] Heater is indispensable in many high latitude countries. Even in a subtropical country like Taiwan, people still need a heater to warm up the temperature in the room. The heaters available in the current market, depending on the way of heating, can be categorized in two different kinds, one is air type heater, and one is radiation type heater. To produce hot air, the air type heater heats the air around the thermal object, and with a fan, to spread the hot air allover the room. Typical air type heaters are ceramic heater, blade-type heater and kerosene radiant heater. The radiation type heater produces heat with the thermal element which is able to emit infrared radiation to warm up the target or the air in a room. The typical radiation type heaters include quartz-tube heater, tungsten lamp heater and halogen lamp heater. [0003] However, the typical disadvantages of air type heater are heating slowly and consuming electrical power, and some of products also exhausting oxygen and making noise. While the radiation type heater although is heating faster but life of product is usually short, and because it's consuming more oxygen, the ventilation become an issue. Other disadvantage of the radiation type heater is that it's emitting red light when operation, which causes disturbance to a sleeper. Normally, no matter which type heater being used, it's hard to match the interior decoration most time so as to become interference in interior design. SUMMARY OF THE INVENTION [0004] In order to cure the disadvantages of traditional heaters described above, the present invention discloses an electrothermal element which is heating faster, without consuming oxygen, with no noise made by fan, occupying less space. The electrothermal element comprises a substrate, which is the main portion of the electrothermal element and can be an object of plate or board; an electrothermal layer, which has electrodes deposited thereon and produces heat and emits infrared radiation when connected electrically; a reflective layer for reflecting infrared radiation from one side of the electrothermal element; and an auxiliary layer which is high thermal conductive, for spreading the heat evenly and converting some portion of thermal energy into infrared radiation so as to enhance the emissivity of infrared radiation of the electrothermal element. [0005] In addition to the advantages described above, the electrothermal element of the present invention also has the following advantages. (1) Heating the target with the infrared radiation, which is much faster than the air type heater. (2) Operating without fan so as to low down the noise. (3) Comparing with the quartz-tube heater, tungsten lamp heater and halogen lamp heater, the operating temperature is low, without the problem of consuming oxygen, life of product is longer, and more safe. (4) Without emitting visible light when in operation to cause interference to the sleeper. (5) The electrothermal element itself is a thin plate or board which can be easily installed in the ceiling or on the wall, and occupying less space than the traditional stand heater. (6) The electrothermal element of the present invention is employing the infrared radiation which enhance the ventilation of blood of human's body, that is, medical effect is excellent, therefore, the present invention is not limited to be used in winter time. (7) The electrothermal element of the present invention radiates infrared radiation more efficient so as to consuming less electrical power to save energy. BRIEF DESCRIPTION OF THE DRAWINGS [0013] The structure and the technical means adopted by the present invention to achieve the above and other objects can be best understood by referring to the following detailed description of the preferred embodiments and the accompanying drawings. [0014] FIG. 1 is a schematic cross-sectional view showing the electrothermal element according to a preferred embodiment of the present invention. [0015] FIG. 2 is a schematic top view showing the electrothermal layer according to a preferred embodiment of the present invention. [0016] FIG. 3 illustrates the transparent reflective layer formed with different doped ZnO films stacked alternatively according to a preferred embodiment of the present invention. [0017] FIG. 4A is a schematic cross-sectional view showing the electrothermal element according to a preferred embodiment of the present invention. [0018] FIG. 4B is a schematic cross-sectional view showing the electrothermal element according to a preferred embodiment of the present invention. [0019] FIG. 4C is a schematic cross-sectional view showing the electrothermal element according to a preferred embodiment of the present invention. [0020] FIG. 4D is a schematic cross-sectional view showing the electrothermal element according to a preferred embodiment of the present invention. [0021] FIG. 5A is a schematic cross-sectional view showing the electrothermal element according to a preferred embodiment of the present invention. [0022] FIG. 5B is a schematic cross-sectional view showing the electrothermal element according to a preferred embodiment of the present invention. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS [0023] Please refer to FIG. 1 , which is the first embodiment of the present invention. In FIG. 1 , the present invention discloses an electrothermal element 1 which comprises a substrate 100 in shape of board or plate, and the substrate 100 can be made of the materials of glass, micro-crystal glass, ceramic or carbon fiber. A reflective layer 400 is deposited on the substrate 100 for reflecting the infrared radiation; the layer can be made of high conductivity metals like gold, silver, copper or aluminum. An auxiliary layer 300 is deposited on the reflective layer 400 for increasing the thermal uniformity of the electrothermal element 1 and also convert the thermal energy of the electrothermal element 1 itself into infrared radiation, which decreases the temperature of the electrothermal element 1 so as to decrease the un-stability of the electrothermal element 1 caused by the temperature, and increase the emissivity of infrared, also stop the aging of the electrothermal element 1 caused by the ion permeation from the substrate (or other layers) to the electrothermal element 1 . The auxiliary layer 300 has the properties of high thermal conductivity and emitting infrared radiation, which can be made of diamond, diamond powder, diamond-like film or diamond-like carbon (DLC) film. An electrothermal layer 200 is deposited on the auxiliary layer 300 which produces heat and emits infrared when being connected electrically, and the electrothermal layer 200 can be made of conductive metal oxides like tin oxide (SnO 2 ), indium tin oxide (ITO) or zinc oxide (ZnO). An electrode 210 is deposited on the electrothermal layer 200 for being connected electrically. The value of electric resistance of the electrothermal layer 200 can be determined by changing the thickness of the layer or by changing the resistivity of the material in production process, or by etching pattern 220 on the electrothermal layer 200 to form an electrical layout. A protective layer 500 is deposited on the electrothermal layer 200 , which can be transparent or non-transparent, for protecting the electrothermal layer 200 from the air, also protecting the object or people from electrical shock. The protective layer 500 can be made of polymeric materials, or the materials used in the auxiliary layer 300 like diamond, diamond powder, diamond-like film or diamond-like carbon (DLC) film so the layer can function as both the protective layer 500 and the auxiliary layer 300 . [0024] As described above, the electrothermal layer 200 produces heat and emits infrared when being connected electrically, the auxiliary layer 300 increases the thermal uniformity and also converts the thermal energy of the electrothermal element itself into infrared radiation, and the reflective layer 400 reflects the infrared radiation efficiently, so, most infrared radiation emits from a first side 11 of the electrothermal element. Because of the advantages described above, the electrothermal element of the present invention is easily installed in the ceiling or on the wall and occupies less space, or can be integrated with the interior design, and the temperature generated by the electrothermal element's is low when in operation, other advantages includes no oxygen consumption and no noise made by fan. [0025] The structure of the second embodiment of the present invention has the similar structure of the first embodiment of the present invention. Please refer to FIG. 1 , which is showing the second embodiment of the present invention, in which the reflective layer 400 is replaced with a transparent and conductive material, such as, transparent conductive films including SnO 2 , ITO, or ZnO, and in this embodiment, ZnO film is used. More particularly, according to the optical theory, the light is reflected at the interface of different materials with different refractive indices. For infrared light, the refractive index of ZnO film is significantly affected by doping levels. A transparent reflective layer 400 (shown in FIG. 3 ) which is transparent in visible light, but high-reflecting in infrared light can be prepared by alternatively stacking different doped ZnO films, such as intrinsic ZnO 410 and doped ZnO 420 . So that, the reflective layer 400 made of different transparent conductive materials, the electrothermal element of the present invention can be an electrothermal element with high transparent index in visible light. [0026] The structure of the third embodiment of the present invention is based on the similar structure of the first embodiment of the present invention, but forming a pattern layer (not shown) on a side 11 of the electrothermal element by printing or sticking to improve the looking of the electrothermal element in order to match the style of interior design when the electrothermal element is installed in the ceiling or on the wall. The pattern layer can be used with thermochromic materials whose color is subjected to the temperature, that is, the temperature change of the electrothermal element will also change the color of the out-looking of the electrothermal element, such feature can be used as an indication of operation or scenario expression. More, the pattern layer can function as a protective layer when the pattern layer is made of insulating material, and when the electrothermal element is transparent with respect to visible light, the pattern layer can be deposited on a second side 12 of the electrothermal element. [0027] The forth embodiment of the present invention is shown in FIG. 4A . In FIG. 4A , the present invention discloses an electrothermal element 1 which comprises a substrate 100 in shape of board or plate, and can be made of the materials of glass, micro-crystal glass, ceramic or carbon fiber. A auxiliary layer 300 is deposited on the substrate 100 for increasing the thermal uniformity of the electrothermal element 1 and also convert the thermal energy of the electrothermal element 1 itself into infrared radiation, which decreases the temperature of the electrothermal element 1 so as to decrease the un-stability of the electrothermal element 1 caused by the temperature, and increase the emissivity of infrared, also stop the aging of the electrothermal element 100 caused by the ion permeation from the substrate (or other layers) to the electrothermal element 1 . The auxiliary layer 300 is high thermal conductive and emitting infrared radiation, which can be made of diamond, diamond powder, diamond-like film or diamond-like carbon (DLC) film. An electrothermal layer 200 is deposited on the auxiliary layer 300 which produces heat and emits infrared when being connected electrically, and the electrothermal layer 200 can be made of conductive metal oxides like SnO 2 , ITO, or ZnO. An electrode 210 is deposited on the electrothermal layer 200 for being connected electrically, and the etching pattern 220 is made on the electrothermal layer 200 for adjusting the value of electric resistance. A reflective layer 400 is deposited on the electrothermal layer 200 for reflecting the infrared radiation; the layer can be made of metals with high conductivity like gold, silver, copper or aluminum. More, the reflective layer 400 can be made to be transparent as illustrated in the second embodiment of the present invention, in which the electrothermal element is a transparent element with respect to visible light. Further, in this forth embodiment, the electrothermal element also comprises a back plate 800 that can be made of the materials of glass, micro-crystal glass, ceramic or carbon fiber. With an adhesive layer 700 , such as PVB or EVA, the back plate 800 can be combined with the substrate 100 . Otherwise, the back plate 800 and the adhesive layer 700 can be replaced together with a protective layer as described in the first embodiment (shown in FIG. 4B ). [0028] The fifth embodiment of the present invention is shown in FIG. 4C , and the structure thereof is similar to that in the forth embodiment of the present invention. In FIG. 4C , the auxiliary layer 300 is deposited on surface 11 of the substrate 100 . Since the auxiliary layer 300 is made of diamond, diamond powder, diamond-like film or diamond-like carbon (DLC) film, which not only has the advantages like lowing down temperature, increasing thermal uniformity of the electrothermal element and increasing the emissivity of infrared, but also can function as a protective layer to protect the surface 11 of the substrate 100 . Moreover, the back plate 800 and the adhesive layer 700 both can be replaced with a protective layer as described in the first embodiment (shown in FIG. 4D ). [0029] The sixth embodiment of the present invention is based on the forth and fifth embodiment of the present invention (shown in FIG. 4A , 4 B, 4 C and 4 D), further forming a pattern layer on the surface 11 of the electrothermal element by printing or sticking to improve the looking of the electrothermal element in order to match the style of interior design when the electrothermal element is installed on the ceiling or on the wall. The pattern layer can be used with thermochromic materials whose color is subjected to the temperature, that is, the temperature change of the electrothermal element will also change the color of the out-looking of the electrothermal element, and such feature can be used as an indication of operation or scenario expression. Moreover, the pattern layer can function as a protective layer when the pattern layer is made of insulating material, and when the electrothermal element is transparent with respect to visible light, the pattern layer can be deposited on the surface 12 of the electrothermal element. [0030] The seventh embodiment of the present invention is shown in FIG. 5A . In FIG. 5 A, the auxiliary layer 300 , the electrothermal layer 200 and a reflecting cover 600 cover on the substrate 100 . The reflective cover 600 can be made of a metal plate or a metal cover, or a cover which of the inner side coated with metals. The function of the reflecting cover 600 is to reflect the infrared and providing protection. The materials of the reflecting cover is made by high conductivity metals, such as gold, silver, copper or aluminum. Moreover, the auxiliary layer is deposited on the surface 11 (as shown in FIG. 5B ) of the electrothermal element. Since the auxiliary layer is made of diamond, diamond powder, diamond-like film or diamond-like carbon (DLC) film, which not only has the advantages like lowing down temperature, increasing thermal uniformity of the electrothermal element and increasing the emissivity of infrared, but also can function as a protective layer to protect the surface of the substrate. Furthermore, a pattern layer can be formed on the surface 11 of the electrothermal element by printing or sticking to improve the out-looking of the electrothermal element in order to match the style of interior design when the electrothermal element is installed on the ceiling or on the wall. The pattern layer can be used with thermochromic materials whose color is subjected to the temperature, that is, the temperature change of the electrothermal element will also change the color of the out-looking of the electrothermal element, and such feature can be used as an indication of operation or scenario expression. [0031] The present invention has been described with some preferred embodiments thereof and it is understood that many changes and modifications in the described embodiments can be carried out without departing from the scope and the spirit of the invention that is intended to be limited only by the appended.

The present invention proposes an electrothermal element which comprises a substrate, a reflective layer, an electrothermal layer, and an auxiliary layer. The electrothermal layer can emit infrared radiation which is reflected by the reflective layer and enable infrared radiation emitted from one side of the electrothermal element. The auxiliary layer can increase the thermal uniformity of the electrothermal element and also convert the thermal energy of the electrothermal element itself into infrared radiation. It not only improves overall infrared emissivity of the electrothermal element but also reduces the temperature of the electrothermal element. This invention offers an effectively and rapidly warming up solution at selected local regions, no oxygen consumption and fan noise problems, and the electrothermal element can replace the traditional ceiling or be mounted directly on the wall, to solve the disadvantage of requiring space of the conventional electric heating device.

STATEMENT OF GOVERNMENT SUPPORT [0001] This invention was made with United States Government support under DARPA/SPO contract BAA00-09. The United States Government may have certain rights in the invention. FIELD OF THE INVENTION [0002] The present invention relates generally to the field of investigational bioinformatics and more particularly to secondary structure defining databases. The present invention further relates to methods for interrogating a database as a source of molecular masses of known bioagents for comparing against the molecular mass of an unknown or selected bioagent to determine either the identity of the selected bioagent, and/or to determine the origin of the selected bioagent. The identification of the bioagent is important for determining a proper course of treatment and/or irradication of the bioagent in such cases as biological warfare. Furthermore, the determination of the geographic origin of a selected bioagent will facilitate the identification of potential criminal identity. The present invention also relates to methods for rapid detection and identification of bioagents from environmental, clinical or other samples. The methods provide for detection and characterization of a unique base composition signature (BCS) from any bioagent, including bacteria and viruses. The unique BCS is used to rapidly identify the bioagent. BACKGROUND OF THE INVENTION [0003] In the United States, hospitals report well over 5 million cases of recognized infectious disease-related illnesses annually. Significantly greater numbers remain undetected, both in the inpatient and community setting, resulting in substantial morbidity and mortality. Critical intervention for infectious disease relies on rapid, sensitive and specific detection of the offending pathogen, and is central to the mission of microbiology laboratories at medical centers. Unfortunately, despite the recognition that outcomes from infectious illnesses are directly associated with time to pathogen recognition, as well as accurate identification of the class and species of microbe, and ability to identify the presence of drug resistance isolates, conventional hospital laboratories often remain encumbered by traditional slow multi-step culture based assays. Other limitations of the conventional laboratory which have become increasingly apparent include: extremely prolonged wait-times for pathogens with long generation time (up to several weeks); requirements for additional testing and wait times for speciation and identification of antimicrobial resistance; diminished test sensitivity for patients who have received antibiotics; and absolute inability to culture certain pathogens in disease states associated with microbial infection. [0004] For more than a decade, molecular testing has been heralded as the diagnostic tool for the new millennium, whose ultimate potential could include forced obsolescence of traditional hospital laboratories. However, despite the fact that significant advances in clinical application of PCR techniques have occurred, the practicing physician still relies principally on standard techniques. A brief discussion of several existing applications of PCR in the hospital-based setting follows. [0005] Generally speaking molecular diagnostics have been championed for identifying organisms that cannot be grown in vitro, or in instances where existing culture techniques are insensitive and/or require prolonged incubation times. PCR-based diagnostics have been successfully developed for a wide variety of microbes. Application to the clinical arena has met with variable success, with only a few assays achieving acceptance and utility. [0006] One of the earliest, and perhaps most widely recognized applications of PCR for clinical practice is in detection of Mycobacterium tuberculosis. Clinical characteristics favoring development of a nonculture-based test for tuberculosis include week to month long delays associated with standard testing, occurrence of drug-resistant isolates and public health imperatives associated with recognition, isolation and treatment. Although frequently used as a diagnostic adjunctive, practical and routine clinical application of PCR remains problematic due to significant inter-laboratory variation in sensitivity, and inadequate specificity for use in low prevalence populations, requiring further development at the technical level. Recent advances in the laboratory suggest that identification of drug resistant isolates by amplification of mutations associated with specific antibiotic resistance (e.g., rpoB gene in rifampin resistant strains) may be forthcoming for clinical use, although widespread application will require extensive clinical validation. [0007] One diagnostic assay, which has gained widespread acceptance, is for C. trachomatis. Conventional detection systems are limiting due to inadequate sensitivity and specificity (direct immunofluoresence or enzyme immunoassay) or the requirement for specialized culture facilities, due to the fastidious characteristics of this microbe. Laboratory development, followed by widespread clinical validation testing in a variety of acute and nonacute care settings have demonstrated excellent sensitivity (90-100%) and specificity (97%) of the PCR assay leading to its commercial development. Proven efficacy of the PCR assay from both genital and urine sampling, have resulted in its application to a variety of clinical setting, most recently including routine screening of patients considered at risk. [0008] While the full potential for PCR diagnostics to provide rapid and critical information to physicians faced with difficult clinical-decisions has yet to be realized, one recently developed assay provides an example of the promise of this evolving technology. Distinguishing life-threatening causes of fever from more benign causes in children is a fundamental clinical dilemma faced by clinicians, particularly when infections of the central nervous system are being considered. Bacterial causes of meningitis can be highly aggressive, but generally cannot be differentiated on a clinical basis from aseptic meningitis, which is a relatively benign condition that can be managed on an outpatient basis. Existing blood culture methods often take several days to turn positive, and are often confounded by poor sensitivity or false-negative findings in patients receiving empiric antimicrobials. Testing and application of a PCR assay for enteroviral meningitis has been found to be highly sensitive. With reporting of results within 1 day, preliminary clinical trials have shown significant reductions in hospital costs, due to decreased duration of hospital stays and reduction in antibiotic therapy. Other viral PCR assays, now routinely available include those for herpes simplex virus, cytomegalovirus, hepatitis and HIV. Each has a demonstrated cost savings role in clinical practice, including detection of otherwise difficult to diagnose infections and newly realized capacity to monitor progression of disease and response to therapy, vital in the management of chronic infectious diseases. [0009] The concept of a universal detection system has been forwarded for identification of bacterial pathogens, and speaks most directly to the possible clinical implications of a broad-based screening tool for clinical use. Exploiting the existence of highly conserved regions of DNA common to all bacterial species in a PCR assay would empower physicians to rapidly identify the presence of bacteremia, which would profoundly impact patient care. Previous empiric decision making could be abandoned in favor of educated practice, allowing appropriate and expeditious decision-making regarding need for antibiotic therapy and hospitalization. [0010] Experimental work using the conserved features of the 16S rRNA common to almost all bacterial species, is an area of active investigation. Hospital test sites have focused on “high yield” clinical settings where expeditious identification of the presence of systemic bacterial infection has immediate high morbidity and mortality consequences. Notable clinical infections have included evaluation of febrile infants at risk for sepsis, detection of bacteremia in febrile neutropenic cancer patients, and examination of critically ill patients in the intensive care unit. While several of these studies have reported promising results (with sensitivity and specificity well over 90%), significant technical difficulties (described below) remain, and have prevented general acceptance of this assay in clinics and hospitals (which remain dependent on standard blood culture methodologies). Even the revolutionary advances of real-time PCR technique, which offers a quantitative more reproducible and technically simpler system remains encumbered by inherent technical limitations of the PCR assay. [0011] The principle shortcomings of applying PCR assays to the clinical setting include: inability to eliminate background DNA contamination; interference with the PCR amplification by substrates present in the reaction; and limited capacity to provide rapid reliable speciation, antibiotic resistance and subtype identification. Some laboratories have recently made progress in identifying and removing inhibitors; however background contamination remains problematic, and methods directed towards eliminating exogenous sources of DNA report significant diminution in assay sensitivity. Finally, while product identification and detailed characterization has been achieved using sequencing techniques, these approaches are laborious and time-intensive thus detracting from its clinical applicability. [0012] Rapid and definitive microbial identification is desirable for a variety of industrial, medical, environmental, quality, and research reasons. Traditionally, the microbiology laboratory has functioned to identify the etiologic agents of infectious diseases through direct examination and culture of specimens. Since the mid-1980s, researchers have repeatedly demonstrated the practical utility of molecular biology techniques, many of which form the basis of clinical diagnostic assays. Some of these techniques include nucleic acid hybridization analysis, restriction enzyme analysis, genetic sequence analysis, and separation and purification of nucleic acids (See, e.g., J. Sambrook, E. F. Fritsch, and T. Maniatis, Molecular Cloning: A Laboratory Manual, 2nd Ed., Cold Spring Harbor Laboratory Press, Cold Spring Harbor, N.Y., 1989). These procedures, in general, are time-consuming and tedious. Another option is the polymerase chain reaction (PCR) or other amplification procedure which amplifies a specific target DNA sequence based on the flanking primers used. Finally, detection and data analysis convert the hybridization event into an analytical result. [0013] Other not yet fully realized applications of PCR for clinical medicine is the identification of infectious causes of disease previously described as idiopathic (e.g. Bartonella henselae in bacillary angiomatosis, and Tropheryma whippellii as the uncultured bacillus associated with Whipple's disease). Further, recent epidemiological studies which suggest a strong association between Chlamydia pneumonia and coronary artery disease, serve as example of the possible widespread, yet undiscovered links between pathogen and host which may ultimately allow for new insights into pathogenesis and novel life sustaining or saving therapeutics. [0014] For the practicing clinician, PCR technology offers a yet unrealized potential for diagnostic omnipotence in the arena of infectious disease. A universal reliable infectious disease detection system would certainly become a fundamental tool in the evolving diagnostic armamentarium of the 21 st century clinician. For front line emergency physicians, or physicians working in disaster settings, a quick universal detection system, would allow for molecular triage and early aggressive targeted therapy. Preliminary clinical studies using species specific probes suggest that implementing rapid testing in acute care setting is feasible. Resources could thus be appropriately applied, and patients with suspected infections could rapidly be risk stratified to the different treatment settings, depending on the pathogen and virulence. Furthermore, links with data management systems, locally regionally and nationally, would allow for effective epidemiological surveillance, with obvious benefits for antibiotic selection and control of disease outbreaks. [0015] For the hospitalists, the ability to speciate and subtype would allow for more precise decision-making regarding antimicrobial agents. Patients who are colonized with highly contagious pathogens could be appropriately isolated on entry into the medical setting without delay. Targeted therapy will diminish development of antibiotic resistance. Furthermore, identification of the genetic basis of antibiotic resistant strains would permit precise pharmacologic intervention. Both physician and patient would benefit with less need for repetitive testing and elimination of wait times for test results. [0016] It is certain that the individual patient will benefit directly from this approach. Patients with unrecognized or difficult to diagnose infections would be identified and treated promptly. There will be reduced need for prolonged inpatient stays, with resultant decreases in iatrogenic events. [0017] Mass spectrometry provides detailed information about the molecules being analyzed, including high mass accuracy. It is also a process that can be easily automated. Low-resolution MS may be unreliable when used to detect some known agents, if their spectral lines are sufficiently weak or sufficiently close to those from other living organisms in the sample. DNA chips with specific probes can only determine the presence or absence of specifically anticipated organisms. Because there are hundreds of thousands of species of benign bacteria, some very similar in sequence to threat organisms, even arrays with 10,000 probes lack the breadth needed to detect a particular organism. [0018] Antibodies face more severe diversity limitations than arrays. If antibodies are designed against highly conserved targets to increase diversity, the false alarm problem will dominate, again because threat organisms are very similar to benign ones. Antibodies are only capable of detecting known agents in relatively uncluttered environments. [0019] Several groups have described detection of PCR products using high resolution electrospray ionization-Fourier transform-ion cyclotron resonance mass spectrometry (ESI-FT-ICR MS). Accurate measurement of exact mass combined with knowledge of the number of at least one nucleotide allowed calculation of the total base composition for PCR duplex products of approximately 100 base pairs. (Aaserud et al., J. Am. Soc. Mass Spec., 1996, 7, 1266-1269; Muddiman et al., Anal. Chem., 1997, 69, 1543-1549; Wunschel et al., Anal. Chem., 1998, 70, 1203-1207; Muddiman et al., Rev. Anal. Chem., 1998, 17, 1-68). Electrospray ionization-Fourier transform-ion cyclotron resistance (ESI-FT-ICR) MS may be used to determine the mass of double-stranded, 500 base-pair PCR products via the average molecular mass (Hurst et al., Rapid Commun. Mass Spec. 1996, 10, 377-382). The use of matrix-assisted laser desorption ionization-time of flight (MALDI-TOF) mass spectrometry for characterization of PCR products has been described. (Muddiman et al., Rapid Commun. Mass Spec., 1999, 13, 1201-1204). However, the degradation of DNAs over about 75 nucleotides observed with MALDI limited the utility of this method. [0020] U.S. Pat. No. 5,849,492 describes a method for retrieval of phylogenetically informative DNA sequences which comprise searching for a highly divergent segment of genomic DNA surrounded by two highly conserved segments, designing the universal primers for PCR amplification of the highly divergent region, amplifying the genomic DNA by PCR technique using universal primers, and then sequencing the gene to determine the identity of the organism. [0021] U.S. Pat. No. 5,965,363 discloses methods for screening nucleic acids for polymorphisms by analyzing amplified target nucleic acids using mass spectrometric techniques and to procedures for improving mass resolution and mass accuracy of these methods. [0022] WO 99/14375 describes methods, PCR primers and kits for use in analyzing preselected DNA tandem nucleotide repeat alleles by mass spectrometry. [0023] WO 98/12355 discloses methods of determining the mass of a target nucleic acid by mass spectrometric analysis, by cleaving the target nucleic acid to reduce its length, making the target single-stranded and using MS to determine the mass of the single-stranded shortened target. Also disclosed are methods of preparing a double-stranded target nucleic acid for MS analysis comprising amplification of the target nucleic acid, binding one of the strands to a solid support, releasing the second strand and then releasing the first strand which is then analyzed by MS. Kits for target nucleic acid preparation are also provided. [0024] PCT WO97/33000 discloses methods for detecting mutations in a target nucleic acid by nonrandomly fragmenting the target into a set of single-stranded nonrandom length fragments and determining their masses by MS. [0025] U.S. Pat. No. 5,605,798 describes a fast and highly accurate mass spectrometer-based process for detecting the presence of a particular nucleic acid in a biological sample for diagnostic purposes. [0026] WO 98/21066 describes processes for determining the sequence of a particular target nucleic acid by mass spectrometry. Processes for detecting a target nucleic acid present in a biological sample by PCR amplification and mass spectrometry detection are disclosed, as are methods for detecting a target nucleic acid in a sample by amplifying the target with primers that contain restriction sites and tags, extending and cleaving the amplified nucleic acid, and detecting the presence of extended product, wherein the presence of a DNA fragment of a mass different from wild-type is indicative of a mutation. Methods of sequencing a nucleic acid via mass spectrometry methods are also described. [0027] WO 97/37041, WO 99/31278 and U.S. Pat. No. 5,547,835 describe methods of sequencing nucleic acids using mass spectrometry. U.S. Pat. Nos. 5,622,824, 5,872,003 and 5,691,141 describe methods, systems and kits for exonuclease-mediated mass spectrometric sequencing. [0028] Thus, there is a need for a method for bioagent detection and identification which is both specific and rapid, and in which no nucleic acid sequencing is required. The present invention addresses this need. SUMMARY OF THE INVENTION [0029] The present invention is directed to method of identifying an unknown bioagent using a database, such as a database stored on, for example, a local computer or perhaps a database accessible over a network or on the internet. This database of molecular masses of known bioagents provides a standard of comparison for determining both identity and geographic origin of the unknown bioagent. The nucleic acid from said bioagent is first contacted with at least one pair of oligonucleotide primers which hybridize to sequences of said nucleic acid that flank a variable nucleic acid sequence of the bioagent. Using PCR technology an amplification product of this variable nucleic acid sequence is made. After standard isolation, the molecular mass of this amplification product is determined using known mass-spec techniques. This molecular mass is compared to the molecular mass of known bioagents within the database, for identifying the unknown bioagent. [0030] This invention is also directed to databases having cell-data positional significance comprising at least a first table that includes a plurality of data-containing cells. The table is organized into at least a first row and a second row, each row having columns which are aligned relative to each other so that inter-row conserved regions are aligned. This alignment facilitates the analysis of regions, which are highly conserved between species. This alignment further provides insight into secondary polymer structure by this alignment. Although this invention is directed to a database where each row describes any polymer, in a preferred embodiment, the polymer is an RNA. Other alignments that operate in the same manner are also contemplated. [0031] Another embodiment of this invention is a method for reconciling the content of two databases such that the content of each is a mirror of the other. [0032] Another embodiment is directed to determining the geographic origin of a bioagent using a database of molecular masses of known bioagents comprising contacting a nucleic acid from the selected bioagent with at least one pair of oligonucleotide primers which hybridize to sequences of the nucleic acid, where the sequences flank a variable nucleic acid sequence of the bioagent. This hybridized region is isolated and amplified through standard PCR techniques known in the art. The molecular mass is determined of this amplified product through any technique known in the art such as, Mass-spectrometry for example. This molecular mass is compared to the molecular masses stored in the database of known bioagents thereby determining a group of probabilistically reasonable geographic origins for the selected bioagent. BRIEF DESCRIPTION OF THE DRAWINGS [0033] FIGS. 1 A- 1 I are representative consensus diagrams that show examples of conserved regions from 16S rRNA (FIGS. 1 A- 1 B), 23S rRNA (3′-half, FIGS. 1 C- 1 D; 5′-half, FIGS. 1 E-F), 23S rRNA Domain I (FIG. 1G), 23S rRNA Domain IV (FIG. 1H) and 16S rRNA Domain III (FIG. 1I) which are suitable for use in the present invention. Lines with arrows are examples of regions to which intelligent primer pairs for PCR are designed. The label for each primer pair represents the starting and ending base number of the amplified region on the consensus diagram. Bases in capital letters are greater than 95% conserved; bases in lower case letters are 90-95% conserved, filled circles are 80-90% conserved; and open circles are less than 80% conserved. The label for each primer pair represents the starting and ending base number of the amplified region on the consensus diagram. [0034] [0034]FIG. 2 shows a typical primer amplified region from the 16S rRNA Domain III shown in FIG. 1C. [0035] [0035]FIG. 3 is a schematic diagram showing conserved regions in RNase P. Bases in capital letters are greater than 90% conserved; bases in lower case letters are 80-90% conserved; filled circles designate bases which are 70-80% conserved; and open circles designate bases that are less than 70% conserved. [0036] [0036]FIG. 4 is a schematic diagram of base composition signature determination using nucleotide analog “tags” to determine base composition signatures. [0037] [0037]FIG. 5 shows the deconvoluted mass spectra of a Bacillus anthracis region with and without the mass tag phosphorothioate A (A*). The two spectra differ in that the measured molecular weight of the mass tag-containing sequence is greater than the unmodified sequence. [0038] [0038]FIG. 6 shows base composition signature (BCS) spectra from PCR products from Staphylococcus aureus ( S. aureus 16S — 1337F) and Bacillus anthracus ( B. anthr. 16S — 1337F), amplified using the same primers. The two strands differ by only two (AT-->CG) substitutions and are clearly distinguished on the basis of their BCS. [0039] [0039]FIG. 7 shows that a single difference between two sequences (A14 in B. anthracis vs. A15 in B. cereus ) can be easily detected using ESI-TOF mass spectrometry. [0040] [0040]FIG. 8 is an ESI-TOF of Bacillus anthracis spore coat protein sspE 56mer plus calibrant. The signals unambiguously identify B. anthracis versus other Bacillus species. [0041] [0041]FIG. 9 is an ESI-TOF of a B. anthracis synthetic 16S — 1228 duplex (reverse and forward strands). The technique easily distinguishes between the forward and reverse strands. [0042] [0042]FIG. 10 is an ESI-FTICR-MS of a synthetic B. anthracis 16S — 1337 46 base pair duplex. [0043] [0043]FIG. 11 is an ESI-TOF-MS of a 56mer oligonucleotide (3 scans) from the B. anthracis saspB gene with an internal mass standard. The internal mass standards are designated by asterisks. [0044] [0044]FIG. 12 is an ESI-TOF-MS of an internal standard with 5 mM TBA-TFA buffer showing that charge stripping with tributylammonium trifluoroacetate reduces the most abundant charge state from [M-8H+]8− to [M-3H+]3−. [0045] [0045]FIG. 13 is a portion of a secondary structure defining database according to one embodiment of the present invention, where two examples of selected sequences are displayed graphically thereunder. [0046] [0046]FIG. 14 is a three dimensional graph demonstrating the grouping of sample molecular weight according to species. [0047] [0047]FIG. 15 is a three dimensional graph demonstrating the grouping of sample molecular weights according to species of virus and mammal infected. [0048] [0048]FIG. 16 is a three dimensional graph demonstrating the grouping of sample molecular weights according to species of virus, and animal-origin of infectious agent. [0049] [0049]FIG. 17 is a figure depicting how the triangulation method of the present invention provides for the identification of an unknown bioagent without prior knowledge of the unknown agent. The use of different primer sets to distinguish and identify the unknown is also depicted as primer sets I, II and III within this figure. A three dimensional graph depicts all of bioagent space ( 170 ), including the unknown bioagent, which after use of primer set I ( 171 ) according to a method according to the present invention further differentiates and classifies bioagents according to major classifications ( 176 ) which, upon further analysis using primer set II ( 172 ) differentiates the unknown agent ( 177 ) from other, known agents ( 173 ) and finally, the use of a third primer set ( 175 ) further specifies subgroups within the family of the unknown ( 174 ). DESCRIPTION OF PREFERRED EMBODIMENTS [0050] The present invention provides a combination of a non-PCR biomass detection mode, preferably high-resolution MS, with PCR-based BCS technology using “intelligent primers” which hybridize to conserved sequence regions of nucleic acids derived from a bioagent and which bracket variable sequence regions that uniquely identify the bioagent. The high-resolution MS technique is used to determine the molecular mass and base composition signature (BCS) of the amplified sequence region. This unique “base composition signature” (BCS) is then input to a maximum-likelihood detection algorithm for matching against a database of base composition signatures in the same amplified region. The present method combines PCR-based amplification technology (which provides specificity) and a molecular mass detection mode (which provides speed and does not require nucleic acid sequencing of the amplified target sequence) for bioagent detection and identification. [0051] The present methods allow extremely rapid and accurate detection and identification of bioagents compared to existing methods. Furthermore, this rapid detection and identification is possible even when sample material is impure. Thus, the method, is useful in a wide variety of fields, including, but not limited to, environmental testing (e.g., detection and discrimination of pathogenic vs. non-pathogenic bacteria in water or other samples), germ warfare (allowing immediate identification of the bioagent and appropriate treatment), pharmacogenetic analysis and medical diagnosis (including cancer diagnosis based on mutations and polymorphisms; drug resistance and susceptibility testing; screening for and/or diagnosis of genetic diseases and conditions; and diagnosis of infectious diseases and conditions). The methods leverage ongoing biomedical research in virulence, pathogenicity, drug resistance and genome sequencing into a method which provides greatly improved sensitivity, specificity and reliability compared to existing methods, with lower rates of false positives. [0052] The present methods can be used, for example, to detect and classify any biological agent, including bacteria, viruses, fungi and toxins. As one example, where the agent is a biological threat, the information obtained is used to determine practical information needed for countermeasures, including toxin genes, pathogenicity islands and antibiotic resistance genes. In addition, the methods can be used to identify natural or deliberate engineering events including chromosome fragment swapping, molecular breeding (gene shuffling) and emerging infectious diseases. [0053] Bacteria have a common set of absolutely required genes. About 250 genes are present in all bacterial species ( Proc. Natl. Acad. Sci. U.S.A., 1996, 93, 10268; Science, 1995, 270, 397), including tiny genomes like Mycoplasma, Ureaplasma and Rickettsia. These genes encode proteins involved in translation, replication, recombination and repair, transcription, nucleotide metabolism, amino acid metabolism, lipid metabolism, energy generation, uptake, secretion and the like. Examples of these proteins are DNA polymerase III beta, elongation factor TU, heat shock protein groEL, RNA polymerase beta, phosphoglycerate kinase, NADH dehydrogenase, DNA ligase, DNA topoisomerase and elongation factor G. Operons can also be targeted using the present method. One example of an operon is the bfp operon from enteropathogenic E. coli. Multiple core chromosomal genes can be used to classify bacteria at a genus or genus species level to determine if an organism has threat potential. The methods can also be used to detect pathogenicity markers (plasmid or chromosomal) and antibiotic resistance genes to confirm the threat potential of an organism and to direct countermeasures. [0054] A theoretically ideal bioagent detector would identify, quantify, and report the complete nucleic acid sequence of every bioagent that reached the sensor. The complete sequence of the nucleic acid component of a pathogen would provide all relevant information about the threat, including its identity and the presence of drug-resistance or pathogenicity markers. This ideal has not yet been achieved. However, the present invention provides a straightforward strategy for obtaining information with the same practical value using base composition signatures (BCS). While the base composition of a gene fragment is not as information-rich as the sequence itself, there is no need to analyze the complete sequence of the gene if the short analyte sequence fragment is properly chosen. A database of reference sequences can be prepared in which each sequence is indexed to a unique base composition signature, so that the presence of the sequence can be inferred with accuracy from the presence of the signature. The advantage of base composition signatures is that they can be quantitatively measured in a massively parallel fashion using multiplex PCR (PCR in which two or more primer pairs amplify target sequences simultaneously) and mass spectrometry. These multiple primer amplified regions uniquely identify most threat and ubiquitous background bacteria and viruses. In addition, cluster-specific primer pairs distinguish important local clusters (e.g., anthracis group). [0055] In the context of this invention, a “bioagent” is any organism, living or dead, or a nucleic acid derived from such an organism. Examples of bioagents include, but are not limited to, cells (including, but not limited to, human clinical samples, bacterial cells and other pathogens), viruses, toxin genes and bioregulating compounds. Samples may be alive or dead or in a vegetative state (for example, vegetative bacteria or spores) and may be encapsulated or bioengineered. [0056] As used herein, a “base composition signature” (BCS) is the exact base composition from selected fragments of nucleic acid sequences that uniquely identifies the target gene and source organism. BCS can be thought of as unique indexes of specific genes. [0057] As used herein, “intelligent primers” are primers which bind to sequence regions which flank an intervening variable region. In a preferred embodiment, these sequence regions which flank the variable region are highly conserved among different species of bioagent. For example, the sequence regions may be highly conserved among all Bacillus species. By the term “highly conserved,” it is meant that the sequence regions exhibit between about 80-100%, more preferably between about 90-100% and most preferably between about 95-100% identity. Examples of intelligent primers which amplify regions of the 16S and 23S rRNA are shown in FIGS. 1 A- 1 I. A typical primer amplified region in 16S rRNA is shown in FIG. 2. The arrows represent primers which bind to highly conserved regions which flank a variable region in 16S rRNA domain III. The amplified region is the stem-loop structure under “1100-1188.” [0058] One main advantage of the detection methods of the present invention is that the primers need not be specific for a particular bacterial species, or even genus, such as Bacillus or Streptomyces. Instead, the primers recognize highly conserved regions across hundreds of bacterial species including, but not limited to, the species described herein. Thus, the same primer pair can be used to identify any desired bacterium because it will bind to the conserved regions which flank a variable region specific to a single species, or common to several bacterial species, allowing nucleic acid amplification of the intervening sequence and determination of its molecular weight and base composition. For example, the 16S — 971-1062, 16S — 1228-1310 and 16S — 1100-1188 regions are 98-99% conserved in about 900 species of bacteria (16S=16S rRNA, numbers indicate nucleotide position). In one embodiment of the present invention, primers used in the present method bind to one or more of these regions or portions thereof. [0059] The present invention provides a combination of a non-PCR biomass detection mode, preferably high-resolution MS, with nucleic acid amplification-based BCS technology using “intelligent primers” which hybridize to conserved regions and which bracket variable regions that uniquely identify the bioagent(s). Although the use of PCR is preferred, other nucleic acid amplification techniques may also be used, including ligase chain reaction (LCR) and strand displacement amplification (SDA). The high-resolution MS technique allows separation of bioagent spectral lines from background spectral lines in highly cluttered environments. The resolved spectral lines are then translated to BCS which are input to a maximum-likelihood detection algorithm matched against spectra for one or more known BCS. Preferably, the bioagent BCS spectrum is matched against one or more databases of BCS from vast numbers of bioagents. Preferably, the matching is done using a maximum-likelihood detection algorithm. [0060] In one embodiment, base composition signatures are quantitatively measured in a massively parallel fashion using the polymerase chain reaction (PCR), preferably multiplex PCR, and mass spectrometric (MS) methods. Sufficient quantities of nucleic acids should be present for detection of bioagents by MS. A wide variety of techniques for preparing large amounts of purified nucleic acids or fragments thereof are well known to those of skill in the art. PCR requires one or more pairs of oligonucleotide primers which bind to regions which flank the target sequence(s) to be amplified. These primers prime synthesis of a different strand of DNA, with synthesis occurring in the direction of one primer towards the other primer. The primers, DNA to be amplified, a thermostable DNA polymerase (e.g. Taq polymerase), the four deoxynucleotide triphosphates, and a buffer are combined to initiate DNA synthesis. The solution is denatured by heating, then cooled to allow annealing of newly added primer, followed by another round of DNA synthesis. This process is typically repeated for about 30 cycles, resulting in amplification of the target sequence. [0061] The “intelligent primers” define the target sequence region to be amplified and analyzed. In one embodiment, the target sequence is a ribosomal RNA (rRNA) gene sequence. With the complete sequences of many of the smallest microbial genomes now available, it is possible to identify a set of genes that defines “minimal life” and identify composition signatures that uniquely identify each gene and organism. Genes that encode core life functions such as DNA replication, transcription, ribosome structure, translation, and transport are distributed broadly in the bacterial genome and are preferred regions for BCS analysis. Ribosomal RNA (rRNA) genes comprise regions that provide useful base composition signatures. Like many genes involved in core life functions, rRNA genes contain sequences that are extraordinarily conserved across bacterial domains interspersed with regions of high variability that are more specific to each species. The variable regions can be utilized to build a database of base composition signatures. The strategy involves creating a structure-based alignment of sequences of the small (16S) and the large (23S) subunits of the rRNA genes. For example, there are currently over 13,000 sequences in the ribosomal RNA database that has been created and maintained by Robin Gutell, University of Texas at Austin, and is publicly available on the Institute for Cellular and Molecular Biology web page on the world wide web of the Internet at, for example, “rna.icmb.utexas.edu/.” There is also a publicly available rRNA database created and maintained by the University of Antwerp, Belgium on the world wide web of the Internet at, for example, “rrna.uia.ac.be.” [0062] These databases have been analyzed to determine regions that are useful as base composition signatures. The characteristics of such regions include: a) between about 80 and 100%, preferably >about 95% identity among species of the particular bioagent of interest, of upstream and downstream nucleotide sequences which serve as sequence amplification primer sites; b) an intervening variable region which exhibits no greater than about 5% identity among species; and c) a separation of between about 30 and 1000 nucleotides, preferably no more than about 50-250 nucleotides, and more preferably no more than about 60-100 nucleotides, between the conserved regions. [0063] Due to their overall conservation, the flanking rRNA primer sequences serve as good “universal” primer binding sites to amplify the region of interest for most, if not all, bacterial species. The intervening region between the sets of primers varies in length and/or composition, and thus provides a unique base composition signature. [0064] It is advantageous to design the “intelligent primers” to be as universal as possible to minimize the number of primers which need to be synthesized, and to allow detection of multiple species using a single pair of primers. These primer pairs can be used to amplify variable regions in these species. Because any variation (due to codon wobble in the 3 rd position) in these conserved regions among species is likely to occur in the third position of a DNA triplet, oligonucleotide primers can be designed such that the nucleotide corresponding to this position is a base which can bind to more than one nucleotide, referred to herein as a “universal base.” For example, under this “wobble” pairing, inosine (I) binds to U, C or A; guanine (G) binds to U or C, and uridine (U) binds to U or C. Other examples of universal bases include nitroindoles such as 5-nitroindole or 3-nitropyrrole (Loakes et al., Nucleosides and Nucleotides, 1995, 14, 1001-1003), the degenerate nucleotides dP or dK (Hill et al.), an acyclic nucleoside analog containing 5-nitroindazole (Van Aerschot et al., Nucleosides and Nucleotides, 1995, 14, 1053-1056) or the purine analog 1-(2-deoxy-β-D-ribofuranosyl)-imidazole-4-carboxamide (Sala et al., Nucl. Acids Res., 1996, 24, 3302-3306). [0065] In another embodiment of the invention, to compensate for the somewhat weaker binding by the “wobble” base, the oligonucleotide primers are designed such that the first and second positions of each triplet are occupied by nucleotide analogs which bind with greater affinity than the unmodified nucleotide. Examples of these analogs include, but are not limited to, 2,6-diaminopurine which binds to thymine, propyne T which binds to adenine and propyne C and phenoxazines, including G-clamp, which binds to G. Propynes are described in U.S. Pat. Nos. 5,645,985, 5,830,653 and 5,484,908, each of which is incorporated herein by reference in its entirety. Phenoxazines are described in U.S. Pat. Nos. 5,502,177, 5,763,588, and 6,005,096, each of which is incorporated herein by reference in its entirety. G-clamps are described in U.S. Pat. Nos. 6,007,992 and 6,028,183, each of which is incorporated herein by reference in its entirety. [0066] Bacterial biological warfare agents capable of being detected by the present methods include, but are not limited to, Bacillus anthracis (anthrax), Yersinia pestis (pneumonic plague), Franciscella tularensis (tularemia), Brucella suis, Brucella abortus, Brucella melitensis (undulant fever), Burkholderia mallei (glanders), Burkholderia pseudomalleii (melioidosis), Salmonella typhi (typhoid fever), Rickettsia typhii (epidemic typhus), Rickettsia prowasekii (endemic typhus) and Coxiella burnetii (Q fever), Rhodobacter capsulatus, Chlamydia pneumoniae, Escherichia coli, Shigella dysenteriae, Shigella flexneri, Bacillus cereus, Clostridium botulinum, Coxiella burnetti, Pseudomonas aeruginosa, Legionella pneumophila, and Vibrio cholerae. [0067] Besides 16S and 23S rRNA, other target regions suitable for use in the present invention for detection of bacteria include, but are not limited to, 5S rRNA and RNase P (FIG. 3). [0068] Biological warfare fungus biowarfare agents include, but are not limited to, coccidioides immitis (Coccidioidomycosis). [0069] Biological warfare toxin genes capable of being detected by the methods of the present invention include, but are not limited to, botulism, T-2 mycotoxins, ricin, staph enterotoxin B, shigatoxin, abrin, aflatoxin, Clostridium perfringens epsilon toxin, conotoxins, diacetoxyscirpenol, tetrodotoxin and saxitoxin. [0070] Biological warfare viral threat agents are mostly RNA viruses (positive-strand and negative-strand), with the exception of smallpox. Every RNA virus is a family of related viruses (quasispecies). These viruses mutate rapidly and the potential for engineered strains (natural or deliberate) is very high. RNA viruses cluster into families that have conserved RNA structural domains on the viral genome (e.g., virion components, accessory proteins) and conserved housekeeping genes that encode core viral proteins including, for single strand positive strand RNA viruses, RNA-dependent RNA polymerase, double stranded RNA helicase, chymotrypsin-like and papain-like proteases and methyltransferases. [0071] Examples of (−)-strand RNA viruses include, but are not limited to, arenaviruses (e.g., sabia virus, lassa fever, Machupo, Argentine hemorrhagic fever, flexal virus), bunyaviruses (e.g., hantavirus, nairovirus, phlebovirus, hantaan virus, Congo-crimean hemorrhagic fever, rift valley fever), and mononegavirales (e.g., filovirus, paramyxovirus, ebola virus, Marburg, equine morbillivirus). [0072] Examples of (+)-strand RNA viruses include, but are not limited to, picornaviruses e.g., coxsackievirus, echovirus, human coxsackievirus A, human echovirus, human enterovirus, human poliovirus, hepatitis A virus, human parechovirus, human rhinovirus), astroviruses (e.g., human astrovirus), calciviruses (e.g., chiba virus, chitta virus, human calcivirus, norwalk virus), nidovirales (e.g., human coronavirus, human torovirus), flaviviruses (e.g., dengue virus 1-4, Japanese encephalitis virus, Kyanasur forest disease virus, Murray Valley encephalitis virus, Rocio virus, St. Louis encephalitis virus, West Nile virus, yellow fever virus, hepatitis c virus) and togaviruses (e.g., Chikugunya virus, Eastern equine encephalitis virus, Mayaro virus, O'nyong-nyong virus, Ross River virus, Venezuelan equine encephalitis virus, Rubella virus, hepatitis E virus). The hepatitis C virus has a 5′-untranslated region of 340 nucleotides, an open reading frame encoding 9 proteins having 3010 amino acids and a 3′-untranslated region of 240 nucleotides. The 5′-UTR and 3′-UTR are 99% conserved in hepatitis C viruses. [0073] In one embodiment, the target gene is an RNA-dependent RNA polymerase or a helicase encoded by (+)-strand RNA viruses, or RNA polymerase from a (−)-strand RNA virus. (+)-strand RNA viruses are double stranded RNA and replicate by RNA-directed RNA synthesis using RNA-dependent RNA polymerase and the positive strand as a template. Helicase unwinds the RNA duplex to allow replication of the single stranded RNA. These viruses include viruses from the family picornaviridae (e.g., poliovirus, coxsackievirus, echovirus), togaviridae (e.g., alphavirus, flavivirus, rubivirus), arenaviridae (e.g., lymphocytic choriomeningitis virus, lassa fever virus), cononaviridae (e.g., human respiratory virus) and Hepatitis A virus. The genes encoding these proteins comprise variable and highly conserved regions which flank the variable regions. [0074] In another embodiment, the detection scheme for the PCR products generated from the bioagent(s) incorporates at least three features. First, the technique simultaneously detects and differentiates multiple (generally about 6-10) PCR products. Second, the technique provides a BCS that uniquely identifies the bioagent from the possible primer sites. Finally, the detection technique is rapid, allowing multiple PCR reactions to be run in parallel. [0075] In one embodiment, the method can be used to detect the presence of antibiotic resistance and/or toxin genes in a bacterial species. For example, Bacillus anthracis comprising a tetracycline resistance plasmid and plasmids encoding one or both anthracis toxins (px01 and/or px02) can be detected by using antibiotic resistance primer sets and toxin gene primer sets. If the B. anthracis is positive for tetracycline resistance, then a different antibiotic, for example quinalone, is used. [0076] Mass spectrometry (MS)-based detection of PCR products provides all of these features with additional advantages. MS is intrinsically a parallel detection scheme without the need for radioactive or fluorescent labels, since every amplification product with a unique base composition is identified by its molecular mass. The current state of the art in mass spectrometry is such that less than femtomole quantities of material can be readily analyzed to afford information about the molecular contents of the sample. An accurate assessment of the molecular mass of the material can be quickly obtained, irrespective of whether the molecular weight of the sample is several hundred, or in excess of one hundred thousand atomic mass units (amu) or Daltons. Intact molecular ions can be generated from amplification products using one of a variety of ionization techniques to convert the sample to gas phase. These ionization methods include, but are not limited to, electrospray ionization (ES), matrix-assisted laser desorption ionization (MALDI) and fast atom bombardment (FAB). For example, MALDI of nucleic acids, along with examples of matrices for use in MALDI of nucleic acids, are described in WO 98/54751 (Genetrace, Inc.). [0077] Upon ionization, several peaks are observed from one sample due to the formation of ions with different charges. Averaging the multiple readings of molecular mass obtained from a single mass spectrum affords an estimate of molecular mass of the bioagent. Electrospray ionization mass spectrometry (ESI-MS) is particularly useful for very high molecular weight polymers such as proteins and nucleic acids having molecular weights greater than 10 kDa, since it yields a distribution of multiply-charged molecules of the sample without causing a significant amount of fragmentation. [0078] The mass detectors used in the methods of the present invention include, but are not limited to, Fourier transform ion cyclotron resonance mass spectrometry (FT-ICR-MS), ion trap, quadrupole, magnetic sector, time of flight (TOF), Q-TOF, and triple quadrupole. [0079] In general, the mass spectrometric techniques which can be used in the present invention include, but are not limited to, tandem mass spectrometry, infrared multiphoton dissociation and pyrolytic gas chromatography mass spectrometry (PGC-MS). In one embodiment of the invention, the bioagent detection system operates continually in bioagent detection mode using pyrolytic GC-MS without PCR for rapid detection of increases in biomass (for example, increases in fecal contamination of drinking water or of germ warfare agents). To achieve minimal latency, a continuous sample stream flows directly into the PGC-MS combustion chamber. When an increase in biomass is detected, a PCR process is automatically initiated. Bioagent presence produces elevated levels of large molecular fragments from, for example, about 100-7,000 Da which are observed in the PGC-MS spectrum. The observed mass spectrum is compared to a threshold level and when levels of biomass are determined to exceed a predetermined threshold, the bioagent classification process described hereinabove (combining PCR and MS, preferably FT-ICR MS) is initiated. Optionally, alarms or other processes (halting ventilation flow, physical isolation) are also initiated by this detected biomass level. [0080] The accurate measurement of molecular mass for large DNAs is limited by the adduction of cations from the PCR reaction to each strand, resolution of the isotopic peaks from natural abundance 13 C and 15 N isotopes, and assignment of the charge state for any ion. The cations are removed by in-line dialysis using a flow-through chip that brings the solution containing the PCR products into contact with a solution containing ammonium acetate in the presence of an electric field gradient orthogonal to the flow. The latter two problems are addressed by operating with a resolving power of >100,000 and by incorporating isotopically depleted nucleotide triphosphates into the DNA. The resolving power of the instrument is also a consideration. At a resolving power of 10,000, the modeled signal from the [M-14H+] 14− charge state of an 84mer PCR product is poorly characterized and assignment of the charge state or exact mass is impossible. At a resolving power of 33,000, the peaks from the individual isotopic components are visible. At a resolving power of 100,000, the isotopic peaks are resolved to the baseline and assignment of the charge state for the ion is straightforward. The [ 13 C, 15 N]-depleted triphosphates are obtained, for example, by growing microorganisms on depleted media and harvesting the nucleotides (Batey et al., Nucl. Acids Res., 1992, 20, 4515-4523). [0081] While mass measurements of intact nucleic acid regions are believed to be adequate to determine most bioagents, tandem mass spectrometry (MS n ) techniques may provide more definitive information pertaining to molecular identity or sequence. Tandem MS involves the coupled use of two or more stages of mass analysis where both the separation and detection steps are based on mass spectrometry. The first stage is used to select an ion or component of a sample from which further structural information is to be obtained. The selected ion is then fragmented using, e.g., blackbody irradiation, infrared multiphoton dissociation, or collisional activation. For example, ions generated by electrospray ionization (ESI) can be fragmented using IR multiphoton dissociation. This activation leads to dissociation of glycosidic bonds and the phosphate backbone, producing two series of fragment ions, called the w-series (having an intact 3′ terminus and a 5′ phosphate following internal cleavage) and the a-Base series (having an intact 5′ terminus and a 3′ furan). [0082] The second stage of mass analysis is then used to detect and measure the mass of these resulting fragments of product ions. Such ion selection followed by fragmentation routines can be performed multiple times so as to essentially completely dissect the molecular sequence of a sample. [0083] If there are two or more targets of similar base composition or mass, or if a single amplification reaction results in a product which has the same mass as two or more bioagent reference standards, they can be distinguished by using mass-modifying “tags.” In this embodiment of the invention, a nucleotide analog or “tag” is incorporated during amplification (e.g., a 5-(trifluoromethyl)deoxythymidine triphosphate) which has a different molecular weight than the unmodified base so as to improve distinction of masses. Such tags are described in, for example, PCT WO97/33000, which is incorporated herein by reference in its entirety. This further limits the number of possible base compositions consistent with any mass. For example, 5-(trifluoromethyl)deoxythymidine triphosphate can be used in place of dTTP in a separate nucleic acid amplification reaction. Measurement of the mass shift between a conventional amplification product and the tagged product is used to quantitate the number of thymidine nucleotides in each of the single strands. Because the strands are complementary, the number of adenosine nucleotides in each strand is also determined. [0084] In another amplification reaction, the number of G and C residues in each strand is determined using, for example, the cytidine analog 5-methylcytosine (5-meC) or propyne C. The combination of the A/T reaction and G/C reaction, followed by molecular weight determination, provides a unique base composition. This method is summarized in FIG. 4 and Table 1. TABLE 1 Total Total Total Base Base base base Double Single mass info info comp. comp. strand strand this this other Top Bottom Mass tag sequence sequence strand strand strand strand strand T* · mass T*ACGT*ACGT* T*ACGT*ACGT* 3x 3T 3A 3T 3A (T* − T) = x AT*GCAT*GCA 2A 2T 2C 2G 2G 2C AT*GCAT*GCA 2x 2T 2A C* · mass TAC*GTAC*GT TAC*GTAC*GT 2x 2C 2G (C* − C) = y ATGC*ATGC*A ATGC*ATGC*A 2x 2C 2G [0085] The mass tag phosphorothioate A (A*) was used to distinguish a Bacillus anthracis cluster. The B. anthracis (A 14 G 9 C 14 T 9 ) had an average MW of 14072.26, and the B. anthracis (A 1 A* 13 G 9 C 14 T 9 ) had an average molecular weight of 14281.11 and the phosphorothioate A had an average molecular weight of +16.06 as determined by ESI-TOF MS. The deconvoluted spectra are shown in FIG. 5. [0086] In another example, assume the measured molecular masses of each strand are 30,000.115 Da and 31,000.115 Da respectively, and the measured number of dT and dA residues are (30,28) and (28,30). If the molecular mass is accurate to 100 ppm, there are 7 possible combinations of dG+dC possible for each strand. However, if the measured molecular mass is accurate to 10 ppm, there are only 2 combinations of dG+dC, and at 1 ppm accuracy there is only one possible base composition for each strand. [0087] Signals from the mass spectrometer may be input to a maximum-likelihood detection and classification algorithm such as is widely used in radar signal processing. The detection processing uses matched filtering of BCS observed in mass-basecount space and allows for detection and subtraction of signatures from known, harmless organisms, and for detection of unknown bioagent threats. Comparison of newly observed bioagents to known bioagents is also possible, for estimation of threat level, by comparing their BCS to those of known organisms and to known forms of pathogenicity enhancement, such as insertion of antibiotic resistance genes or toxin genes. [0088] Processing may end with a Bayesian classifier using log likelihood ratios developed from the observed signals and average background levels. The program emphasizes performance predictions culminating in probability-of-detection versus probability-of-false-alarm plots for conditions involving complex backgrounds of naturally occurring organisms and environmental contaminants. Matched filters consist of a priori expectations of signal values given the set of primers used for each of the bioagents. A genomic sequence database (e.g. GenBank) is used to define the mass basecount matched filters. The database contains known threat agents and benign background organisms. The latter is used to estimate and subtract the signature produced by the background organisms. A maximum likelihood detection of known background organisms is implemented using matched filters and a running-sum estimate of the noise covariance. Background signal strengths are estimated and used along with the matched filters to form signatures which are then subtracted. the maximum likelihood process is applied to this “cleaned up” data in a similar manner employing matched filters for the organisms and a running-sum estimate of the noise-covariance for the cleaned up data. [0089] In one embodiment, a strategy to “triangulate” each organism by measuring signals from multiple core genes is used to reduce false negative and false positive signals, and enable reconstruction of the origin or hybrid or otherwise engineered bioagents. After identification of multiple core genes, alignments are created from nucleic acid sequence databases. The alignments are then analyzed for regions of conservation and variation, and potential primer binding sites flanking variable regions are identified. Next, amplification target regions for signature analysis are selected which distinguishes organisms based on specific genomic differences (i.e., base composition). For example, detection of signatures for the three part toxin genes typical of B. anthracis (Bowen et al., J. Appl. Microbiol., 1999, 87, 270-278) in the absence of the expected signatures from the B. anthracis genome would suggest a genetic engineering event. [0090] The present method can also be used to detect single nucleotide polymorphisms (SNPs), or multiple nucleotide polymorphisms, rapidly and accurately. A SNP is defined as a single base pair site in the genome that is different from one individual to another. The difference can be expressed either as a deletion, an insertion or a substitution, and is frequently linked to a disease state. Because they occur every 100-1000 base pairs, SNPs are the most frequently bound type of genetic marker in the human genome. [0091] For example, sickle cell anemia results from an A-T transition, which encodes a valine rather than a glutamic acid residue. Oligonucleotide primers may be designed such that they bind to sequences which flank a SNP site, followed by nucleotide amplification and mass determination of the amplified product. Because the molecular masses of the resulting product from an individual who does not have sickle cell anemia is different from that of the product from an individual who has the disease, the method can be used to distinguish the two individuals. Thus, the method can be used to detect any known SNP in an individual and thus diagnose or determine increased susceptibility to a disease or condition. [0092] In one embodiment, blood is drawn from an individual and peripheral blood mononuclear cells (PBMC) are isolated and simultaneously tested, preferably in a high-throughput screening method, for one or more SNPs using appropriate primers based on the known sequences which flank the SNP region. The National Center for Biotechnology Information maintains a publicly available database of SNPs on the world wide web of the Internet at, for example, “ncbi.nlm.nih.gov/SNP/.” [0093] The method of the present invention can also be used for blood typing. The gene encoding A, B or O blood type can differ by four single nucleotide polymorphisms. If the gene contains the sequence CGTGGTGACCCTT (SEQ ID NO:1), antigen A results. If the gene contains the sequence CGTCGTCACCGCTA (SEQ ID NO:2) antigen B results. If the gene contains the sequence CGTGGT-ACCCCTT (SEQ ID NO:3), blood group O results (“-” indicates a deletion). These sequences can be distinguished by designing a single primer pair which flanks these regions, followed by amplification and mass determination. [0094] While the present invention has been described with specificity in accordance with certain of its preferred embodiments, the following examples serve only to illustrate the invention and are not intended to limit the same. EXAMPLES Example 1 Nucleic Acid Isolation and PCR [0095] In one embodiment, nucleic acid is isolated from the organisms and amplified by PCR using standard methods prior to BCS determination by mass spectrometry. Nucleic acid is isolated, for example, by detergent lysis of bacterial cells, centrifugation and ethanol precipitation. Nucleic acid isolation methods are described in, for example, Current Protocols in Molecular Biology (Ausubel et al.) and Molecular Cloning; A Laboratory Manual (Sambrook et al.). The nucleic acid is then amplified using standard methodology, such as PCR, with primers which bind to conserved regions of the nucleic acid which contain an intervening variable sequence as described below. Example 2 Mass Spectrometry [0096] FTICR Instrumentation: The FTICR instrument is based on a 7 tesla actively shielded superconducting magnet and modified Bruker Daltonics Apex II 70e ion optics and vacuum chamber. The spectrometer is interfaced to a LEAP PAL autosampler and a custom fluidics control system for high throughput screening applications. Samples are analyzed directly from 96-well or 384-well microtiter plates at a rate of about 1 sample/minute. The Bruker data-acquisition platform is supplemented with a lab-built ancillary NT datastation which controls the autosampler and contains an arbitrary waveform generator capable of generating complex rf-excite waveforms (frequency sweeps, filtered noise, stored waveform inverse Fourier transform (SWIFT), etc.) for sophisticated tandem MS experiments. For oligonucleotides in the 20-30-mer regime typical performance characteristics include mass resolving power in excess of 100,000 (FWHM), low ppm mass measurement errors, and an operable m/z range between 50 and 5000 m/z. [0097] Modified ESI Source: In sample-limited analyses, analyte solutions are delivered at 150 nL/minute to a 30 mm i.d. fused-silica ESI emitter mounted on a 3-D micromanipulator. The ESI ion optics consist of a heated metal capillary, an rf-only hexapole, a skimmer cone, and an auxiliary gate electrode. The 6.2 cm rf-only hexapole is comprised of 1 mm diameter rods and is operated at a voltage of 380 Vpp at a frequency of 5 MHz. A lab-built electromechanical shutter can be employed to prevent the electrospray plume from entering the inlet capillary unless triggered to the “open” position via a TTL pulse from the data station. When in the “closed” position, a stable electrospray plume is maintained between the ESI emitter and the face of the shutter. The back face of the shutter arm contains an elastomeric seal which can be positioned to form a vacuum seal with the inlet capillary. When the seal is removed, a 1 mm gap between the shutter blade and the capillary inlet allows constant pressure in the external ion reservoir regardless of whether the shutter is in the open or closed position. When the shutter is triggered, a “time slice” of ions is allowed to enter the inlet capillary and is subsequently accumulated in the external ion reservoir. The rapid response time of the ion shutter (<25 ms) provides reproducible, user defined intervals during which ions can be injected into and accumulated in the external ion reservoir. [0098] Apparatus for Infrared Multiphoton Dissociation: A 25 watt CW CO 2 laser operating at 10.6 μm has been interfaced to the spectrometer to enable infrared multiphoton dissociation (IRMPD) for oligonucleotide sequencing and other tandem MS applications. An aluminum optical bench is positioned approximately 1.5 m from the actively shielded superconducting magnet such that the laser beam is aligned with the central axis of the magnet. Using standard IR-compatible mirrors and kinematic mirror mounts, the unfocused 3 mm laser beam is aligned to traverse directly through the 3.5 mm holes in the trapping electrodes of the FTICR trapped ion cell and longitudinally traverse the hexapole region of the external ion guide finally impinging on the skimmer cone. This scheme allows IRMPD to be conducted in an m/z selective manner in the trapped ion cell (e.g. following a SWIFT isolation of the species of interest), or in a broadband mode in the high pressure region of the external ion reservoir where collisions with neutral molecules stabilize IRMPD-generated metastable fragment ions resulting in increased fragment ion yield and sequence coverage. Example 3 Identification of Bioagents [0099] Table 2 shows a small cross section of a database of calculated molecular masses for over 9 primer sets and approximately 30 organisms. The primer sets were derived from rRNA alignment. Examples of regions from rRNA consensus alignments are shown in FIGS. 1 A- 1 C. Lines with arrows are examples of regions to which intelligent primer pairs for PCR are designed. The primer pairs are >95% conserved in the bacterial sequence database (currently over 10,000 organisms). The intervening regions are variable in length and/or composition, thus providing the base composition “signature” (BCS) for each organism. Primer pairs were chosen so the total length of the amplified region is less than about 80-90 nucleotides. The label for each primer pair represents the starting and ending base number of the amplified region on the consensus diagram. [0100] Included in the short bacterial database cross-section in Table 2 are many well known pathogens/biowarfare agents (shown in bold/red typeface) such as Bacillus anthracis or Yersinia pestis as well as some of the bacterial organisms found commonly in the natural environment such as Streptomyces. Even closely related organisms can be distinguished from each other by the appropriate choice of primers. For instance, two low G+C organisms, Bacillus anthracis and Staph aureus, can be distinguished from each other by using the primer pair defined by 16S — 1337 or 23S — 855 (ΔM of 4 Da). TABLE 2 Cross Section Of A Database Of Calculated Molecular Masses 1 Primer Regions→ Bug Name 16S_971 16S_1100 16S_1337 16S_1294 16S_1228 23S_1021 23S_855 23S_193 23S_115 Acinetobacter calcoaceticus 55619.1 55004 28446.7 35854.9 51295.4 30299 42654 39557.5 54999 Bacillus anthracis 55005 54388 28448 35238 51296 30295 42651 39560 56850 Bacillus cereus 55622.1 54387.9 28447.6 35854.9 51296.4 30295 42651 39560.5 56850.3 Bordetella bronchiseptica 56857.3 51300.4 28446.7 35857.9 51307.4 30299 42653 39559.5 51920.5 Borrelia burgdorferi 56231.2 55621.1 28440.7 35852.9 51295.4 30297 42029.9 38941.4 52524.6 Brucella abortus 58098 55011 28448 35854 50683 Campylobacter jejuni 58088.5 54386.9 29061.8 35856.9 50674.3 30294 42032.9 39558.5 45732.5 Chlamydia pnuemoniae 55000 55007 29063 35855 50676 30295 42036 38941 56230 Clostridlum botullinum 55006 53767 28445 35855 51291 30300 42656 39562 54999 Clostridium difficile 56855.3 54386.9 28444.7 35853.9 51296.4 30294 41417.8 39556.5 55612.2 Enterococcus faecalis 55620.1 54387.9 28447.6 35858.9 51296.4 30297 42652 39559.5 56849.3 Escherichia coli 55622 55009 28445 35857 51301 30301 42656 39562 54999 Francisella tularensis 53769 54385 28445 35856 51298 Haemophilus influenzae 55620.1 55006 28444.7 35855.9 51298.4 30298 42656 39560.5 55613.1 Klebsiella pneumoniae 55622.1 55008 28442.7 35856.9 51297.4 30300 42655 39562.5 55000 Legionella pneumophila 55618 55626 28446 35857 51303 Mycobacterium avium 54390.9 55631.1 29084.8 35858.9 51915.5 30298 42656 38942.4 56241.2 Mycobacterium leprae 54389.9 55629.1 29084.8 35860.9 51917.5 30298 42656 39559.5 56240.2 Mycobacterium tuberculosis 54390.9 55629.1 29064.8 35860.9 51301.4 30299 42656 39560.5 56243.2 Mycoplasma genitalium 53143.7 45115.4 29061.8 35854.9 50671.3 30294 43264.1 39558.5 56842.4 Mycoplasma pneumoniae 53143.7 45118.4 29061.8 35854.9 50673.3 30294 43264.1 39559.5 56843.4 Neisseria gonorrhoeae 55627.1 54389.9 28445.7 35855.9 51302.4 30300 42649 39561.5 55000 Pseudomonas aeruginosa 55623 55010 28443 35858 51301 30298 43272 39558 55619 Rickettsia prowazekii 58093 55621 28448 35853 50677 30293 42650 39559 53139 Rickettsia rickettsii 58094 55623 28448 35853 50679 30293 42648 39559 53755 Salmonella typhimurium 55622 55005 28445 35857 51301 30301 42658 Shigella dysenteriae 55623 55009 28444 35857 51301 Staphylococcus aureus 56854.3 54386.9 28443.7 35852.9 51294.4 30298 42655 39559.5 57466.4 Streptomyces 54389.9 59341.6 29063.8 35858.9 51300.4 39563.5 56864.3 Treponema pallidum 56245.2 55631.1 28445.7 35851.9 51297.4 30299 42034.9 38939.4 57473.4 Vibrio cholerae 55625 55626 28443 35857 52536 29063 30303 35241 50675 Vibrio parahaemolyticus 54384.9 55626.1 28444.7 34620.7 50064.2 Yersinia pestis 55620 55626 28443 35857 51299 [0101] [0101]FIG. 6 shows the use of ESI-FT-ICR MS for measurement of exact mass. The spectra from 46mer PCR products originating at position 1337 of the 16S rRNA from S. aureus (upper) and B. anthracis (lower) are shown. These data are from the region of the spectrum containing signals from the [M-8H+] 8− charge states of the respective 5′-3′ strands. The two strands differ by two (AT→CG) substitutions, and have measured masses of 14206.396 and 14208.373±0.010 Da, respectively. The possible base compositions derived from the masses of the forward and reverse strands for the B. anthracis products are listed in Table 3. TABLE 3 Possible base composition for B. anthracis products Calc. Mass Error Base Comp. 14208.2935 0.079520 A1 G17 C10 T18 14208.3160 0.056980 A1 G20 C15 T10 14208.3386 0.034440 A1 G23 C20 T2 14208.3074 0.065560 A6 G11 C3 T26 14208.3300 0.043020 A6 G14 C8 T18 14208.3525 0.020480 A6 G17 C13 T10 14208.3751 0.002060 A6 G20 C18 T2 14208.3439 0.029060 A11 G8 C1 T26 14208.3665 0.006520 A11 G11 C6 T18 14208.3890 0.016020 A11 G14 C11 T10 14208.4116 0.038560 A11 G17 C16 T2 14208.4030 0.029980 A16 G8 C4 T18 14208.4255 0.052520 A16 G11 C9 T10 14208.4481 0.075060 A16 G14 C14 T2 14208.4395 0.066480 A21 G5 C2 T18 14208.4620 0.089020 A21 G8 C7 T10 14079.2624 0.080600 A0 G14 C13 T19 14079.2849 0.058060 A0 G17 C18 T11 14079.3075 0.035520 A0 G20 C23 T3 14079.2538 0.089180 A5 G5 C1 T35 14079.2764 0.066640 A5 G8 C6 T27 14079.2989 0.044100 A5 G11 C11 T19 14079.3214 0.021560 A5 G14 C16 T11 14079.3440 0.000980 A5 G17 C21 T3 14079.3129 0.030140 A10 G5 C4 T27 14079.3354 0.007600 A10 G8 C9 T19 14079.3579 0.014940 A10 G11 C14 T11 14079.3805 0.037480 A10 G14 C19 T3 14079.3494 0.006360 A15 G2 C2 T27 14079.3719 0.028900 A15 G5 C7 T19 14079.3944 0.051440 A15 G5 C12 T11 14079.4170 0.073980 A15 G11 C17 T3 14079.4084 0.065400 A20 G2 C5 T19 14079.4309 0.087940 A20 G5 C10 T13 [0102] Among the 16 compositions for the forward strand and the 18 compositions for the reverse strand that were calculated, only one pair (shown in bold) are complementary, corresponding to the actual base compositions of the B. anthracis PCR products. Example 4 BCS of Region from Bacillus anthracis and Bacillus cereus [0103] A conserved Bacillus region from B. anthracis (A 14 G 9 C 14 T 9 ) and B. cereus (A 15 G 9 C 13 T 9 ) having a C to A base change was synthesized and subjected to ESI-TOF MS. The results are shown in FIG. 7 in which the two regions are clearly distinguished using the method of the present invention (MW=14072.26 vs. 14096.29). Example 5 Identification of Additional Bioagents [0104] In other examples of the present invention, the pathogen Vibrio cholera can be distinguished from Vibrio parahemolyticus with ΔM>600 Da using one of three 16S primer sets shown in Table 2 (16S — 971, 16S — 1228 or 16S — 1294) as shown in Table 4. The two mycoplasma species in the list ( M. genitalium and M. pneumoniae ) can also be distinguished from each other, as can the three mycobacteriae. While the direct mass measurements of amplified products can identify and distinguish a large number of organisms, measurement of the base composition signature provides dramatically enhanced resolving power for closely related organisms. In cases such as Bacillus anthracis and Bacillus cereus that are virtually indistinguishable from each other based solely on mass differences, compositional analysis or fragmentation patterns are used to resolve the differences. The single base difference between the two organisms yields different fragmentation patterns, and despite the presence of the ambiguous/unidentified base N at position 20 in B. anthracis, the two organisms can be identified. [0105] Tables 4a-b show examples of primer pairs from Table 1 which distinguish pathogens from background. TABLE 4a Organism name 23S_855 16S_1337 23S_1021 Bacillus anthracis 42650.98 28447.65 30294.98 Staphylococcus aureus 42654.97 28443.67 30297.96 [0106] [0106] TABLE 4b Organism name 16S_971 16S_1294 16S_1228 Vibrio cholerae 55625.09 35856.87 52535.59 Vibrio parahaemolyticus 54384.91 34620.67 50064.19 [0107] Table 4 shows the expected molecular weight and base composition of region 16S — 1100-1188 in Mycobacterium avium and Streptomyces sp. TABLE 5 Molecular Region Organism name Length weight Base comp. 16S_1100-1188 Mycobacterium avium 82 25624.1728 A 16 G 32 C 18 T 16 16S_1100-1188 Streptomyces sp. 96 29904.871 A 17 G 38 C 27 T 14 [0108] Table 5 shows base composition (single strand) results for 16S — 1100-1188 primer amplification reactions different species of bacteria. Species which are repeated in the table (e.g., Clostridium botulinum ) are different strains which have different base compositions in the 16S — 1100-1188 region. TABLE 6 Organism name Base comp. Organism name Base comp. Mycobactenum avium A 16 G 32 C 18 T 16 Vibrio cholerae A 23 G 30 C 21 T 16 Streptomyces sp. A 17 G 38 C 27 T 14 Aeromonas hydrophila A 23 G 31 C 21 T 15 Ureaplasma urealyticum A 18 G 30 C 17 T 17 Aeromonas salmonicida A 23 G 31 C 21 T 15 Streptomyces sp. A 19 G 36 C 24 T 18 Mycoplasma genitalium A 24 G 19 C 12 T 18 Mycobacterium leprae A 20 G 32 C 22 T 16 Clostridium botulinum A 24 G 25 C 18 T 20 M. tuberculosis A 20 G 33 C 21 T 16 Bordetella bronchiseptica A 24 G 26 C 19 T 14 Nocardia asteroides A 20 G 33 C 21 T 16 Francisella tularensis A 24 G 26 C 19 T 19 Fusobacterium A 21 G 26 C 22 T 18 Bacillus anthracis A 24 G 26 C 20 T 18 necroforum Listeria monocytogenes A 21 G 27 C 19 T 19 Campylobacter jejuni A 24 G 26 C 20 T 18 Clostridium botulinum A 21 G 27 C 19 T 21 Staphylococcus aureus A 24 G 26 C 20 T 18 Neisseria gonorrhoeae A 21 G 28 C 21 T 18 Helicobacter pylori A 24 G 26 C 20 T 19 Bartonella quintana A 21 G 30 C 22 T 16 Helicobacter pylori A 24 G 26 C 21 T 18 Enterococcus faecalis A 22 G 27 C 20 T 19 Moraxella catarrhalis A 24 G 26 C 23 T 16 Bacillus megaterium A 22 G 28 C 20 T 18 Haemophilus influenzae Rd A 24 G 28 C 20 T 17 Bacillus subtilis A 22 G 28 C 21 T 17 Chlamydia trachomatis A 24 G 28 C 21 T 16 Pseudomonas aeruginosa A 22 G 29 C 23 T 15 Chlamydophila pneumoniae A 24 G 28 C 21 T 16 Legionella pneumophila A 22 G 32 C 20 T 16 C. pneumonia AR39 A 24 G 28 C 21 T 16 Mycoplasma pneumoniae A 23 G 20 C 14 T 16 Pseudomonas putida A 24 G 29 C 21 T 16 Clostridium botulinum A 23 G 26 C 20 T 19 Proteus vulgaris A 24 G 30 C 21 T 15 Enterococcus faecium A 23 G 26 C 21 T 18 Yersinia pestis A 24 G 30 C 21 T 15 Acinetobacter calcoaceti A 23 G 26 C 21 T 19 Yersinia pseudotuberculos A 24 G 30 C 21 T 15 Leptospira borgpeterseni A 23 G 26 C 24 T 15 Clostridium botulinum A 25 G 24 C 18 T 21 Leptospira interrogans A 23 G 26 C 24 T 15 Clostridium tetani A 25 G 25 C 18 T 20 Clostridium perfringens A 23 G 27 C 19 T 19 Francisella tularensis A 25 G 25 C 19 T 19 Bacillus anthracis A 23 G 27 C 20 T 18 Acinetobacter calcoacetic A 25 G 26 C 20 T 19 Bacillus cereus A 23 G 27 C 20 T 18 Bacteriodes fragilis A 25 G 27 C 16 T 22 Bacillus thuringiensis A 23 G 27 C 20 T 18 Chlamydophila psittaci A 25 G 27 C 21 T 16 Aeromonas hydrophila A 23 G 29 C 21 T 16 Borrelia burgdorfeni A 25 G 29 C 17 T 19 Escherichia coli A 23 G 29 C 21 T 16 Streptobacillus monilifor A 26 G 26 C 20 T 16 Pseudomonas putida A 23 G 29 C 21 T 17 Rickettsia prowazekii A 26 G 28 C 18 T 18 Escherichia coil A 23 G 29 C 11 T 15 Rickettsia rickettsii A 26 G 28 C 20 T 16 Shigella dysenteriae A 23 G 29 C 22 T 15 Mycoplasma mycoides A 28 G 23 C 16 T 20 [0109] The same organism having different base compositions are different strains. Groups of organisms which are highlighted or in italics have the same base compositions in the amplified region. Some of these organisms can be distinguished using multiple primers. For example, Bacillus anthracis can be distinguished from Bacillus cereus and Bacillus thuringiensis using the primer 16S — 971-1062 (Table 6). Other primer pairs which produce unique base composition signatures are shown in Table 6 (bold). Clusters containing very similar threat and ubiquitous non-threat organisms (e.g. anthracis cluster) are distinguished at high resolution with focused sets of primer pairs. The known biowarfare agents in Table 6 are Bacillus anthracis, Yersinia pestis, Francisella tularensis and Rickettsia prowazekii . TABLE 7 Organism 16S_971-1062 16S_1228-1310 16S_1100-1188 Aeromonas hydrophila A 21 G 29 C 22 T 20 A 22 G 27 C 21 T 13 A 23 G 31 C 21 T 15 Aeromonas salmonicida A 21 G 29 C 22 T 20 A 22 G 27 C 21 T 13 A 23 G 31 C 21 T 15 Bacillus anthracis A 21 G 27 C 22 T 22 A 24 G 22 C 19 T 18 A 23 G 27 C 20 T 18 Bacillus cereus A 22 G 27 C 21 T 22 A 24 G 22 C 19 T 18 A 23 G 27 C 20 T 18 Bacillus thuringiensis A 22 G 27 C 21 T 22 A 24 G 22 C 19 T 18 A 23 G 27 C 20 T 18 Chlamydia trachomatis A 22 G 26 C 20 T 23 A 24 G 23 C 19 T 16 A 24 G 28 C 21 T 16 Chlamydia pneumoniae AR39 A 26 G 23 C 20 T 22 A 26 G 22 C 16 T 18 A 24 G 28 C 21 T 16 Leptospira borgpetersenii A 22 G 26 C 20 T 21 A 22 G 25 C 21 T 15 A 23 G 26 C 24 T 15 Leptospira interrogans A 22 G 26 C 20 T 21 A 22 G 25 C 21 T 15 A 23 G 26 C 24 T 15 Mycoplasma genitalium A 28 G 23 C 15 T 22 A 30 G 18 C 15 T 19 A 24 G 19 C 12 T 18 Mycoplasma pneumoniae A 28 G 23 C 15 T 22 A 27 G 19 C 16 T 20 A 23 G 20 C 14 T 16 Escherichia coli A 22 G 28 C 20 T 22 A 24 G 25 C 21 T 13 A 23 G 29 C 22 T 15 Shigella dysenteriae A 22 G 28 C 21 T 21 A 24 G 25 C 21 T 13 A 23 G 29 C 22 T 15 Proteus vulgaris A 23 G 26 C 22 T 21 A 26 G 24 C 19 T 14 A 24 G 30 C 21 T 15 Yersinia pestis A 24 G 25 C 21 T 22 A 25 G 24 C 20 T 14 A 24 G 30 C 21 T 15 Yersinia pseudotuberculosis A 24 G 25 C 21 T 22 A 25 G 24 C 20 T 14 A 24 G 30 C 21 T 15 Francisella tularensis A 20 G 25 C 21 T 23 A 23 G 26 C 17 T 17 A 24 G 26 C 19 T 19 Rickettsia prowazekii A 21 G 26 C 24 T 25 A 24 G 23 C 16 T 19 A 26 G 28 C 18 T 18 Rickettsia rickettsii A 21 G 26 C 25 T 24 A 24 G 24 C 17 T 17 A 26 G 28 C 20 T 16 [0110] The sequence of B. anthracis and B. cereus in region 16S — 971 is shown below. Shown in bold is the single base difference between the two species which can be detected using the methods of the present invention. B. anthracis has an ambiguous base at position 20. B. anthracis _16S_971 (SEQ ID NO:4) GCGAAGAACCUUACCAGGU N UUGACAUCCUCUGA C AACCCUAGAGAUAGG GCUUCUCCUUCGGGAGCAGAGUGACAGGUGGUGCAUGGUU B. cereus _16S_971 (SEQ ID NO:5) GCGAAG AA CCUUACCAGGUCUUGACAUCCUCUG AA AACCCUAGAGAUAGG GCUUCUCCUUCGGGAGCAGAGUGACAGGUGGUGCAUGGUU Example 6 ESI-TOF MS of sspE 56-mer Plus Calibrant [0111] The mass measurement accuracy that can be obtained using an internal mass standard in the ESI-MS study of PCR products is shown in FIG. 8. The mass standard was a 20-mer phosphorothioate oligonucleotide added to a solution containing a 56-mer PCR product from the B. anthracis spore coat protein sspE. The mass of the expected PCR product distinguishes B. anthracis from other species of Bacillus such as B. thuringiensis and B. cereus. Example 7 B. anthracis ESI-TOF Synthetic 16S — 1228 Duplex [0112] An ESI-TOF MS spectrum was obtained from an aqueous solution containing 5 μM each of synthetic analogs of the expected forward and reverse PCR products from the nucleotide 1228 region of the B. anthracis 16S rRNA gene. The results (FIG. 9) show that the molecular weights of the forward and reverse strands can be accurately determined and easily distinguish the two strands. The [M-21H + ] 21− and [M-20H + ] 20− charge states are shown. Example 8 ESI-FTICR-MS of Synthetic B. anthracis 16S — 1337 46 Base Pair Duplex [0113] An ESI-FTICR-MS spectrum was obtained from an aqueous solution containing 5 μM each of synthetic analogs of the expected forward and reverse PCR products from the nucleotide 1337 region of the B. anthracis 16S rRNA gene. The results (FIG. 10) show that the molecular weights of the strands can be distinguished by this method. The [M-16H + ] 16− through [M-10H + ] 10− charge states are shown. The insert highlights the resolution that can be realized on the FTICR-MS instrument, which allows the charge state of the ion to be determined from the mass difference between peaks differing by a single 13C substitution. Example 9 ESI-TOF MS of 56-mer Oligonucleotide from saspB Gene of B. anthracis with Internal Mass Standard [0114] ESI-TOF MS spectra were obtained on a synthetic 56-mer oligonucleotide (5 μM)from the saspB gene of B. anthracis containing an internal mass standard at an ESI of 1.7 μL/min as a function of sample consumption. The results (FIG. 11) show that the signal to noise is improved as more scans are summed, and that the standard and the product are visible after only 100 scans. Example 10 ESI-TOF MS of an Internal Standard with Tributylammonium (TBA)-Trifluoroacetate (TFA) Buffer [0115] An ESI-TOF-MS spectrum of a 20-mer phosphorothioate mass standard was obtained following addition of 5 mM TBA-TFA buffer to the solution. This buffer strips charge from the oligonucleotide and shifts the most abundant charge state from [M-8H + ] 8− to [M-3H + ] 3− (FIG. 12). Example 11 Master Database Comparison [0116] The molecular masses obtained through Examples 1-10 are compared to molecular masses of known bioagents stored in a master database to obtain a high probability matching molecular mass. Example 12 Master Data Base Interrogation over the Internet [0117] The same procedure as in Example 11 is followed except that the local computer did not store the Master database. The Master database is interrogated over an internet connection, searching for a molecular mass match. Example 13 Master Database Updating [0118] The same procedure as in example 11 is followed except the local computer is connected to the internet and has the ability to store a master database locally. The local computer system periodically, or at the user's discretion, interrogates the Master database, synchronizing the local master database with the global Master database. This provides the current molecular mass information to both the local database as well as to the global Master database. This further provides more of a globalized knowledge base. Example 14 Global Database Updating [0119] The same procedure as in example 13 is followed except there are numerous such local stations throughout the world. The synchronization of each database adds to the diversity of information and diversity of the molecular masses of known bioagents. [0120] Various modifications of the invention, in addition to those described herein, will be apparent to those skilled in the art from the foregoing description. Such modifications are also intended to fall within the scope of the appended claims. 1 7 1 90 RNA Bacillus anthracis misc_feature (20)..(20) N = A, U, G or C 1 gcgaagaacc uuaccaggun uugacauccu cugacaaccc uagagauagg gcuucuccuu 60 cgggagcaga gugacaggug gugcaugguu 90 2 90 RNA Bacillus cereus 2 gcgaagaacc uuaccagguc uugacauccu cugaaaaccc uagagauagg gcuucuccuu 60 cgggagcaga gugacaggug gugcaugguu 90 3 1542 RNA Artificial Sequence misc_feature 16S rRNA consensus sequence 3 nnnnnnnaga guuugaucnu ggcucagnnn gaacgcuggc ggnnngcnun anacaugcaa 60 gucgancgnn nnnnnnnnnn nnnnnnnnnn nnnnnnnnnn agnggcnnac gggugaguaa 120 nncnunnnna nnunccnnnn nnnnnggnan annnnnnnga aannnnnnnu aauaccnnau 180 nnnnnnnnnn nnnnaaagnn nnnnnnnnnn nnnnnnnnnn nnnnnngann nnnnnnngnn 240 nnaunagnun guuggunngg uaanggcnna ccaagncnnn gannnnuagc ngnncugaga 300 ggnngnncng ccacanuggn acugaganac ggnccanacu ccuacgggag gcagcagunn 360 ggaaunuunn ncaauggnng naanncugan nnagcnannc cgcgugnnng anganggnnu 420 nnngnungua aannncunun nnnnnngang annnnnnnnn nnnnnnnnnn nnnnnnnnnu 480 gacnnuannn nnnnannaag nnncggcnaa cuncgugcca gcagccgcgg uaauacgnag 540 gnngcnagcg uunnncggan unanugggcg uaaagngnnn gnaggnggnn nnnnnngunn 600 nnngunaaan nnnnnngcun aacnnnnnnn nnncnnnnnn nacnnnnnnn cungagnnnn 660 nnagnggnnn nnngaauunn nnguguagng gugnaauncg naganaunng nangaanacc 720 nnungcgaag gcnnnnnncu ggnnnnnnac ugacncunan nnncgaaagc nugggnagcn 780 aacaggauua gauacccugg uaguccangc nnuaaacgnu gnnnnnunnn ngnnngnnnn 840 nnnnnnnnnn nnnnnnnnna nnnaacgnnn uaannnnncc gccuggggag uacgnncgca 900 agnnunaaac ucaaangaau ugacggggnc cngcacaagc ngnggagnau guggnuuaau 960 ucgangnnac gcgnanaacc uuaccnnnnn uugacaunnn nnnnnnnnnn nnganannnn 1020 nnnnnnnnnn nnnnnnnnnn nnnacaggug nugcauggnu gucgucagcu cgugnnguga 1080 gnuguugggu uaagucccgn aacgagcgca acccnnnnnn nnnguuncna ncnnnnnnnn 1140 ngngnacucn nnnnnnacug ccnnngnnaa nnnggaggaa ggnggggang acgucaanuc 1200 nucaugnccc uuangnnnng ggcuncacac nuncuacaau ggnnnnnaca nngngnngcn 1260 annnngnnan nnnnagcnaa ncnnnnaaan nnnnucnnag uncggaungn nnncugcaac 1320 ucgnnnncnu gaagnnggan ucgcuaguaa ucgnnnauca gnangnnncg gugaauacgu 1380 ucncgggncu uguacacacc gcccgucann ncangnnagn nnnnnnnncc nnaagnnnnn 1440 nnnnnnncnn nnnngnnnnn nnnnncnang gnnnnnnnnn nganugggnn naagucguaa 1500 caagguancc nuannngaan nugnggnugg aucaccuccu un 1542 4 2904 RNA Artificial Sequence misc_feature 23S rRNA consensus sequence 4 nnnnaagnnn nnaagngnnn nngguggaug ccunggcnnn nnnagncgan gaaggangnn 60 nnnnncnncn nnanncnnng gnnagnngnn nnnnnncnnn nnanccnnng nunuccgaau 120 ggggnaaccc nnnnnnnnnn nnnnnnnnan nnnnnnnnnn nnnnnnnnnn nnnnnnngnn 180 nacnnnnnga anugaaacau cunaguannn nnaggaanag aaannaannn ngauuncnnn 240 nguagnggcg agcgaannng nannagncnn nnnnnnnnnn nnnnnnnnnn nnnannngaa 300 nnnnnuggna agnnnnnnnn nannngguna nannccngua nnnnaaannn nnnnnnnnnn 360 nnnnnnnnnn aguannncnn nncncgngnn annnngunng aannngnnnn gaccannnnn 420 naagncuaaa uacunnnnnn ngaccnauag ngnannagua cngugangga aaggngaaaa 480 gnacccnnnn nangggagug aaanagnncc ugaaaccnnn nncnuanaan nngunnnagn 540 nnnnnnnnnn nnnuganngc gunccuuuug nannaugnnn cngnganuun nnnunnnnng 600 cnagnuuaan nnnnnnnngn agncgnagng aaancgagun nnaanngngc gnnnagunnn 660 nngnnnnaga cncgaancnn ngugancuan nnaugnncag gnugaagnnn nnguaanann 720 nnnuggaggn ccgaacnnnn nnnnguugaa aannnnnngg augannugug nnungnggng 780 aaanncnaan cnaacnnngn nauagcuggu ucucnncgaa annnnuuuag gnnnngcnun 840 nnnnnnnnnn nnnnggnggu agagcacugn nnnnnnnnng gnnnnnnnnn nnnnuacnna 900 nnnnnnnnaa acuncgaaun ccnnnnnnnn nnnnnnnngn agnnanncnn ngngngnuaa 960 nnuncnnngu nnanagggna acancccaga ncnncnnnua aggncccnaa nnnnnnnnua 1020 aguggnaaan gangugnnnn nncnnanaca nnnaggangu uggcuuagaa gcagccancn 1080 uunaaagann gcguaanagc ucacunnucn agnnnnnnng cgcngannau nuancgggnc 1140 uaannnnnnn nccgaannnn nngnnnnnnn nnnnnnnnnn nnnnngguag nngagcgunn 1200 nnnnnnnnnn ngaagnnnnn nngnnannnn nnnuggannn nnnnnnagug ngnaugnngn 1260 naunaguanc gannnnnnnn gugananncn nnnncnccgn annncnaagg nuuccnnnnn 1320 nangnunnuc nnnnnngggu nagucgnnnc cuaagnngag ncnganangn nuagnngaug 1380 gnnannnggu nnauauuccn nnacnnnnnn nnnnnnnnnn nnnnngacgn nnnnngnnnn 1440 nnnnnnnnnn nnnnggnnnn nnnnnnnnnn nnnnnnnnnn nnnnnnnnnn nnnnnnnnnn 1500 nnnnnnnnnn nnnnnnnnnn nnnnnnnnnn nnnnnnnnnn nnnnnnnnnn nnnnnnnnnn 1560 nnnncnngaa aannnnnnnn nnnnnnnnnn nnnnnnnnnc guaccnnaaa ccgacacagg 1620 ungnnnngnn gagnanncnn aggngnnngn nnnaannnnn nnnaaggaac unngcaaanu 1680 nnnnccguan cuucggnana aggnnnncnn nnnnnnnnnn nnnnnnnnnn nnnnnnnnnn 1740 nnnnnnnnng nnnnannnan nngnnnnnnn cnacuguuua nnaaaaacac agnncnnugc 1800 naanncgnaa gnnganguau anggnnugac nccugcccng ugcnngaagg uuaanngnnn 1860 nnnnnngnnn nngnnnnnnn nnnnannnaa gcccnnguna acggcggnng uaacuauaac 1920 nnuccuaagg uagcgaaauu ccuugucggg uaaguuccga ccngcacgaa nggngnaang 1980 annnnnnnnc ugucucnnnn nnnnncncng ngaanuunna nunnnnguna agaugcnnnn 2040 uncncgcnnn nngacggaaa gaccccnngn ancuuuacun nannnunnna nugnnnnnnn 2100 nnnnnnnnug unnagnauag gunggagncn nngannnnnn nncgnnagnn nnnnnggagn 2160 cnnnnnugnn auacnacncu nnnnnnnnnn nnnnucuaac nnnnnnnnnn nancnnnnnn 2220 nnngacanug nnngnngggn aguuunacug gggcggunnc cuccnaaann guaacggagg 2280 ngnncnaagg unnncunann nnggnnggnn aucnnnnnnn nagunnaann gnanaagnnn 2340 gcnunacugn nagnnnnacn nnncgagcag nnncgaaagn nggnnnuagu gauccggngg 2400 unnnnnnugg aagngccnuc gcucaacgga uaaaagnuac ncnggggaua acaggcunau 2460 nnnncccaag aguncanauc gacggnnnng uuuggcaccu cgaugucggc ucnucncauc 2520 cuggggcugn agnngguccc aagggunngg cuguucgccn nuuaaagngg nacgngagcu 2580 ggguunanaa cgucgugaga caguungguc ccuaucngnn gngngngnnn gannnuugan 2640 nngnnnugnn cnuaguacga gaggaccggn nngnacnnan cncuggugnn ncnguugunn 2700 ngccannngc anngcngnnu agcuannunn ggnnnngaua anngcugaan gcaucuaagn 2760 nngaancnnn cnnnnagann agnnnucncn nnnnnnnnnn nnnnnnnnna gnnncnnnnn 2820 agannannnn gungauaggn nngnnnugna agnnnngnna nnnnunnagn nnacnnnuac 2880 uaaunnnncn nnnnncuunn nnnn 2904 5 13 DNA Artificial Sequence misc_feature Primer 5 cgtggtgacc ctt 13 6 14 DNA Artificial Sequence misc_feature Primer 6 cgtcgtcacc gcta 14 7 13 DNA Artificial Sequence misc_feature Primer 7 cgtggtaccc ctt 13

The present invention relates generally to the field of investigational bioinformatics and more particularly to secondary structure defining databases. The present invention further relates to methods for interrogating a database as a source of molecular masses of known bioagents for comparing against the molecular mass of an unknown or selected bioagent to determine either the identity of the selected bioagent, and/or to determine the origin of the selected bioagent. The identification of the bioagent is important for determining a proper course of treatment and/or irradication of the bioagent in such cases as biological warfare. Furthermore, the determination of the geographic origin of a selected bioagent will facilitate the identification of potential criminal identity.

BACKGROUND OF THE INVENTION [0001] 1. Field of the Invention [0002] The present invention relates to a geared motor including a gear train and the like inside of a case, a manufacturing method for the geared motor, and a damper device used in a cold air passage of a refrigerator or the like. [0003] 2. Description of Related Art [0004] In a damper device which is used in a cold air passage of a refrigerator or the like, for example, a structure has been proposed that a baffle is driven by a baffle drive mechanism including a motor and a gear train to open and close an opening part formed in a frame (Patent Literature 1). In the damper device, the baffle drive mechanism is accommodated in a case to structure a geared motor. In the geared motor and the damper device, motor lead wires are arranged within the case and, in Patent Literature 1, the motor lead wires are connected with a connector near the motor. CITATION LIST Patent Literature [0005] [Patent Literature 1] Japanese Patent Laid-Open No. 2010-159902 SUMMARY OF THE INVENTION [0006] However, in the structure described in Patent Literature 1, a power feeding position to the motor is limited and thus it is inconvenient to use. On the other hand, when a connector is disposed at a position separated from the motor, the motor lead wires may be contacted with the gear train. [0007] In view of the problem described above, an objective of the present invention is to provide a geared motor and a damper device capable of leading around motor lead wires to a position separated from a motor inside the case. Means to Solve the Problems [0008] To solve the above mentioned problem, the present invention provides a geared motor, wherein directions perpendicular to each other are referred to as an “X” direction and a “Y” direction, and a direction perpendicular to the “X” direction and the “Y” direction is referred to as a “Z” direction, the geared motor including a case, which is a bottomed case including a case body part which opens toward one side in the “X” direction and a bottom plate part located on the other side in the “X”, direction with respect to the case body part, a cover which covers an opening of the case body part of the case on the one side in the “X” direction, a motor which is disposed inside the case, a plurality of motor lead wires having flexibility which are connected with the motor, and a gear train disposed on one side in the “Y” direction with respect to the motor in the case in a manner that turning center axial lines of the gear train are directed in the “X” direction. The case includes a motor lead wire passage whose depth direction is the “X” direction and which is extended in the “Y” direction so as to pass a position interposed in the “Z” direction between the turning center axial lines of a gear included in the gear train and the case body part, and the plurality of the motor lead wires are extended from the other side to the one side in the “Y” direction through the motor lead wire passage in a state that at least parts of the plurality of the motor lead wires are bent in the motor lead wire passage from a connecting position of the motor lead wires and the motor so as to be located on a side of the bottom plate part. [0009] In the present invention, the case is provided with a motor lead wire passage which is extended in the “Y” direction so as to pass a position interposed in the “Z” direction between the turning center axial lines of the gears included in the gear train and the case body part. Therefore, a plurality of the motor lead wires can be extended in the “Y” direction in a state that the positions in the “Z” direction of the motor lead wires are restricted. Further, the plurality of the motor lead wires is extended in the “Y” direction through the motor lead wire passage in a state that at least parts of the motor lead wires are bent to a side of the bottom plate part from the connecting positions in the motor lead wire passage. Therefore, the motor lead wires can be led around to a position separated in the “Y” direction from the motor in a state that the motor lead wires and the gear train are not contacted with each other. Accordingly, flexibility for a power feeding position to the motor can be enhanced. [0010] In the present invention, it may be structured that the motor lead wire passage is formed in a slit shape which opens toward the one side in the “X” direction. According to this structure, the motor lead wires can be fitted to the motor lead wire passage from the one side in the “X” direction. [0011] In the present invention, it is preferable that parts of the plurality of the motor lead wires are overlapped with the gears on the other side in the “X” direction in the motor lead wire passage when viewed in the “X” direction. According to this structure, a dimension in the “Z” direction of the geared motor can be reduced. [0012] In this case, it is preferable that the motor lead wire passage is structured by the bottom plate part on the other side in the “X” direction. According to this structure, parts of the motor lead wires can be bent to a position contacted with the bottom plate part or to a position close to the bottom plate part. Therefore, when viewed in the “X” direction, even in a case that a plurality of the motor lead wires and the gears are overlapped with each other, the motor lead wires and the gears are hard to be contacted with each other. [0013] In the present invention, it is preferable that the motor lead wire is provided with a first portion which is extended from the connecting position, a second portion which is bent from an end part of the first portion to the other side in the “X” direction and is extended in the motor lead wire passage, a third portion which is bent from an end part on the other side in the “X” direction of the second portion to the one side in the “Y” direction and is extended in the motor lead wire passage, a fourth portion which is bent from an end part on the one side in the “Y” direction of the third portion to the one side in the “X” direction and is extended in the motor lead wire passage, and a fifth portion which is extended from an end part on the one side in the “X” direction of the fourth portion, and at least a part of the third portion is overlapped with the gears on the other side in the “X” direction when viewed in the “X” direction. According to this structure, when viewed in the “X” direction, even in a case that a plurality of the motor lead wires and the gears are overlapped with each other, the motor lead wires and the gears are hard to be contacted with each other. [0014] In this case, it is preferable that the fifth portion is extended in the “Z” direction and held by the case. According to this structure, even if the fifth portion is pulled, the third portion is hard to be floated to the one side in the “X” direction. Therefore, when viewed in the “X” direction, even in a case that a plurality of the motor lead wires and the gears are overlapped with each other, the motor lead wires and the gears are hard to be contacted with each other. [0015] In the present invention, it is preferable that the cover is provided with a cover side projection at a position displaced from the gears in the “Y” direction so as to protrude toward the other side in the “X” direction to prevent displacement of the third portion to the one side in the “X” direction. According to this structure, displacement to the one side in the “X” direction of the third portion can be prevented and thus, when viewed in the “X” direction, even in a case that the third portion and a gears are overlapped with each other, the motor lead wires and the gears are hard to be contacted with each other. [0016] In the present invention, it is preferable that the plurality of the motor lead wires is juxtaposed in the “X” direction and extended in the “Y” direction in the motor lead wire passage. According to this structure, a region occupied by the plurality of the motor lead wires is narrow in the “Z” direction and thus even when the plurality of the motor lead wires are disposed at a position near the turning center axial lines of the gears, the motor lead wires and the gears are hard to be contacted with each other. [0017] In the present invention, it is preferable that the case is provided with a lead wire receiving part which is obliquely inclined from the “X” direction to the “Z” direction in the motor lead wire passage, and the third portions of the plurality of the motor lead wires are juxtaposed in the “X” direction in an obliquely inclined state in the “Z” direction along the lead wire receiving part. According to this structure, a region occupied in the “X” direction by the plurality of the motor lead wires is narrow and thus, even when the third portions and the gears are overlapped with each other when viewed in the “X” direction, the motor lead wires and the gears are hard to be contacted with each other. [0018] In the present invention, it is preferable that the plurality of the motor lead wires are formed in a flat cable in which the plurality of the motor lead wires are connected with each other in a juxtaposed state in the “X” direction in the motor lead wire passage. According to this structure, a plurality of the motor lead wires is easily accommodated in the motor lead wire passage. [0019] In the present invention, it is preferable that lengths of part or all of the plurality of the motor lead wires are different from each other, the motor includes a plurality of terminals whose distances from an outlet of the motor lead wire passage are different from each other, and part or all of the plurality of the terminals are connected with the motor lead wires having shorter lengths as located nearer to the outlet of the motor lead wire passage. According to this structure, the terminal farther from the outlet of the motor lead wire passage is connected with a motor lead wire having a longer length. Therefore, a motor lead wire having a needlessly long length is prevented from being connected with a terminal which is near to the outlet of the motor lead wire passage and thus slack of the lead wire does not become excessively large. Further, a shorter motor lead wire can be avoided from being connected with a terminal which is far from the outlet of the motor lead wire passage and thus the motor lead wire is avoided from being pulled. Therefore, the motor lead wires can be easily led around from the motor lead wire passage to the terminals. [0020] In the present invention, it is preferable that, in the motor lead wire passage, the plurality of the motor lead wires is located in the motor lead wire passage closer to the one side in the “X” direction as the terminal connected with the motor lead wire is nearer to the outlet of the motor lead wire passage. According to this structure, the nearer to the outlet of the motor lead wire passage a terminal is, a motor lead wire overlapped on an upper side in the motor lead wire passage is connected with the nearer terminal. Therefore, a plurality of the motor lead wires can be led around orderly. [0021] The present invention is effectively applied to a case that the case body part is formed in a rectangular shape whose long side is extended in the “Y” direction when viewed in the “X” direction. [0022] In the present invention, it is preferable that the motor is held by the case and the gear train is turnably supported by the bottom plate part of the case. According to this structure, it is sufficient to fit the cover after the motor and the gear train are assembled in the case and thus assembling work is easily performed. [0023] The present invention provides a manufacturing method for the above-mentioned geared motor including a bending tendency giving process in which the plurality of the motor lead wires are given with a bending tendency in a bent shape, and a assembling process in which the plurality of the motor lead wires which have been given with the bending tendency is assembled in the case. According to this structure, when the motor lead wires are to be assembled in the case, bending work is not required to perform at the spot. Therefore, work for assembling the motor lead wires is easily performed and a working hour can be shortened. [0024] The geared motor to which the present invention is applied can be used in a damper device. The damper device includes a frame which is provided on the one side in the “X” direction with respect to the case and is formed with an opening part opening in the “Z” direction, and a baffle configured to open and close the opening part, and the baffle is driven by the geared motor. [0025] In the present invention, it is preferable that the cover is integrally structured with the frame. According to this structure, in comparison with a case that the cover and the frame are separately structured from each other, assembling efficiency is improved and the number of components can be reduced. [0026] In the present invention, it is preferable that the damper device includes a heater which is fixed to the frame around the opening part and a heater lead wire having flexibility which is electrically connected with the heater. The heater lead wire is extended to the other side in the “X” direction of the cover through a heater lead wire passage formed in the cover and is led out from the case to an outer side together with the plurality of the motor lead wires, and at least one of the cover and the case is provided with a lead wire support part which supports a midway portion of the heater lead wire from an outer side. According to this structure, the heater lead wire is held at the time of assembling of the damper device. Therefore, troublesome labor is not required for fixing the heater lead wire. [0027] In this case, it is preferable that the cover is provided with a cover side lead wire support part as the lead wire support part which supports midway portions of the plurality of the motor lead wires and the heater lead wire from the outer side between an outer face of the case and the cover side lead wire support part. Effects of the Invention [0028] In the present invention, the case is provided with a motor lead wire passage which is extended in the “Y” direction so as to pass a position interposed in the “Z” direction between a turning center axial lines of gears included in a gear train and the case body part. Therefore, a plurality of the motor lead wires can be extended in the “Y” direction in a state that the positions in the “Z” direction of the motor lead wires are restricted. Further, the plurality of the motor lead wires is extended in the “Y” direction through the motor lead wire passage in a state that at least parts of the motor lead wires are bent to a side of the bottom plate part with respect to the connecting positions with the motor in the motor lead wire passage. Therefore, the motor lead wires can be led around to a position separated in the “Y” direction from the motor in a state that the motor lead wires and the gear train are not contacted with each other. Accordingly, flexibility for a power feeding position to the motor can be enhanced. BRIEF DESCRIPTION OF THE DRAWINGS [0029] FIGS. 1( a ), 1( b ) and 1( c ) are perspective views showing a damper device to which the present invention is applied and which is viewed from a side where a baffle is disposed. [0030] FIGS. 2( a ) and 2( b ) are perspective views showing a damper device to which the present invention is applied and which is viewed from an opposite side to a side where a baffle is disposed. [0031] FIGS. 3( a ), 3( b ) and 3( c ) are explanatory views showing a leading-around structure of motor lead wires in a geared motor and a damper device to which the present invention is applied. [0032] FIGS. 4( a ), 4( b ) and 4( c ) are explanatory views showing an inside of a case in a geared motor and a damper device to which the present invention is applied. [0033] FIG. 5 is an explanatory view showing a motor lead wire passage which is formed in a case in a geared motor and a damper device to which the present invention is applied. [0034] FIG. 6 is a perspective outward appearance view showing a motor. [0035] FIGS. 7( a ), 7( b ) and 7( c ) are explanatory views schematically showing motor lead wires in a modified embodiment. [0036] FIGS. 8( a ), 8( b ) and 8( c ) are explanatory views showing a jig for performing connecting work of motor lead wires and a motor. DESCRIPTION OF EMBODIMENTS [0037] A damper device for a refrigerator to which the present invention is applied will be described below with reference to the accompanying drawings. In the following descriptions, a turning center axial line of a baffle 4 is indicated by the “L”, a direction along the turning center axial line “L” is referred to as the “X” direction, a direction that an opening part 210 faces is referred to as the “Z” direction, and a direction perpendicular to the “X” direction and the “Z” direction is referred to as the “Y” direction. Further, the “X1” is one side in the “X” direction, the “X2” is the other side in the “X” direction, the “Y1” is one side in the “Y” direction, the “Y2” is the other side in the “Y” direction, the “Z1” is one side in the “Z” direction, and the “Z2” is the other side in the “Z” direction. (Entire Structure) [0038] FIGS. 1( a ), 1( b ) and 1( c ) are perspective views showing a damper device 1 to which the present invention is applied and which is viewed from a side where a baffle 4 is disposed. FIG. 1( a ) is a perspective view showing an entire damper device 1 , FIG. 1( b ) is an exploded perspective view showing the damper device 1 which is disassembled into a frame 2 and a case 3 , and FIG. 1( c ) is an exploded perspective view showing the damper device 1 in which a baffle 4 is detached. FIGS. 2( a ) and 2( b ) are perspective views showing the damper device 1 to which the present invention is applied and which is viewed from an opposite side to a side where the baffle 4 is disposed. FIG. 2( a ) is a perspective view showing the entire damper device 1 and FIG. 2( b ) is an exploded perspective view showing the damper device 1 which is disassembled into the frame 2 and the case 3 . FIGS. 1( a ), 1( b ) and 1( c ) and FIGS. 2( a ) and 2( b ) show a state that the baffle 4 closes an opening part 210 . Further, in FIG. 1( b ) , heater lead wires 8 are shown so as to be held on a side of the case 3 for convenience but, when the case 3 is to be connected with a cover 20 , the heater lead wires 8 are held on the side of the cover 20 . [0039] As shown in FIGS. 1( a ), 1( b ) and 1( c ) and FIGS. 2( a ) and 2( b ) , the damper device 1 in this embodiment includes a flame 2 formed with a rectangular opening part 210 which opens in the “Z” direction, a baffle 4 for opening and closing the opening part 210 of the frame 2 , a baffle drive mechanism 6 which is disposed in a drive chamber 60 adjacent to the opening part 210 and the baffle 4 through a cover 20 on the other side “X2” in the “X” direction, and a case 3 which is fitted to the cover 20 from the other side “X2” in the “X” direction to section the drive chamber 60 between the cover 20 and the case 3 . In this embodiment, the cover 20 and the case 3 are connected by a hook mechanism 39 . In this state, the case 3 , the cover 20 and the baffle drive mechanism 6 structure a geared motor 10 . The cover 20 is formed with shaft parts 29 protruded to the other side “X2” in the “X” direction to position the case 3 with respect to the cover 20 and a cover side projection 28 which presses motor lead wires 7 described below. In this embodiment, the frame 2 , the cover 20 and the case 3 are made of resin. [0040] The frame 2 is provided with a rectangular end plate part 21 in which the opening part 210 is formed, and a body part 22 in a rectangular tube shape which is protruded to the other side “Z2” in the “Z” direction from an outer side edge of the end plate part 21 . In this embodiment, the cover 20 is structured as a portion of the body part 22 located on the other side “X2” in the “X” direction and is integrally formed with the frame 2 . [0041] A seal plate part 23 in a rectangular tube shape which is protruded toward a side where the baffle 4 is located is formed at an edge of the opening part 210 in the end plate part 21 . The baffle 4 is abutted with the seal plate part 23 to set the opening part 210 in a closed state. The baffle 4 includes an opening/closing plate 40 provided with a flat plate part 41 whose size is larger than the opening part 210 , and a sheet-shaped elastic member 49 made of foamed polyurethane or the like which is stuck on a face of the opening/closing plate 40 on a side of the opening part 210 . The elastic member 49 is abutted with a portion surrounding the opening part 210 (seal plate part 23 ) to close the opening part 210 . [0042] The baffle 4 is supported by the frame 2 so as to be turnable around the turning center axial line “L” extended in the “X” direction and the baffle drive mechanism 6 turns the baffle 4 around the turning center axial line “L” to open and close the opening part 210 . [0043] In this embodiment, a face of the end plate part 21 of the frame 2 on a side where the baffle 4 is located is attached with a heater 9 so as to surround the opening part 210 (surround the seal plate part 23 ). In this embodiment, the heater 9 is formed in a sheet shape. [0044] The damper device 1 is disposed on an inner side of a duct which structures a cold air passage. In this embodiment, cold air flows through the opening part 210 from an opposite side to a side where the baffle 4 is disposed with respect to the opening part 210 . Alternatively, cold air may flow through the opening part 210 from a side where the baffle 4 is disposed with respect to the opening part 210 . (Structure of Baffle Drive Mechanism 6 ) [0045] In the damper device 1 and the geared motor 10 in this embodiment, the baffle drive mechanism 6 includes a motor 61 disposed on an inner side of the case 3 and a gear train 65 structured to transmit rotation of the motor 61 to the baffle 4 on one side “Y1” in the “Y” direction with respect to the motor 61 . In this embodiment, the motor 61 is a stepping motor. The case 3 is a bottomed case which is provided with a bottom plate part 31 located on the other side “X2” in the “X” direction and a rectangular tube shaped case body part 32 protruded to a side of the cover 20 (one side “X1” in the “X” direction) from the bottom plate part 31 . The case body part 32 is opened toward the one side “X1” in the “X” direction. The case body part 32 is provided with side plate parts 321 and 322 facing in the “Z” direction and side plate parts 323 and 324 facing in the “Y” direction. The bottom plate part 31 and the case body part 32 are, when viewed in the “X” direction, formed in a quadrangular shape whose long sides are extended in the “Y” direction and short sides are extended in the “Z” direction. The motor 61 is held between the frame 2 and the cover 20 by the case 3 . [0046] The gear train 65 includes a first gear 66 having a large diameter gear engaged with a motor pinion, a second gear 67 having a large diameter gear engaged with a small diameter gear of the first gear 66 , and a drive gear 68 having a large diameter gear engaged with a small diameter gear of the second gear 67 . Therefore, in the gear train 65 , the first gear 66 , the second gear 67 and the drive gear 68 structure a reduction gear train. Further, the gear train 65 includes a sector gear 69 which is engaged with the drive gear 68 and is driven by the drive gear 68 . The sector gear 69 is a final gear (output gear) located at the last stage of the gear train 65 and is connected with the baffle 4 . In this embodiment, all the gears of the gear train 65 (first gear 66 , second gear 67 , drive gear 68 and sector gear 69 ) are turnably supported by the bottom plate part 31 of the case 3 with their turning center axial lines directing in the “X” direction. [0047] In the damper device 1 structured as described above, the motor 61 is connected with totaled four (4) motor lead wires 7 . The motor lead wires 7 are, inside the drive chamber 60 , led around from a connecting position 79 with the motor 61 to one side “Y1” in the “Y” direction through the other side “Z2” in the “Z” direction and then extended toward the other side “X2” in the “X” direction and led out to the outside of the case 3 . [0048] Further, in the damper device 1 , the heater 9 is connected with totaled two (2) heater lead wires 8 . The heater lead wires 8 are connected with the heater 9 on one side “X1” in the “X” direction (side where the heater 9 is disposed) with respect to the cover 20 and then led around to the other side “X2” in the “X” direction with respect to the cover 20 through a heater lead wire passage 25 formed in the cover 20 , and then extended toward the other side “X2” in the “X” direction and led out to the outside of the case 3 . The heater lead wires 8 are extended toward the other side “X2” in the “X” direction juxtaposed with the motor lead wires 7 and led out to the outside of the case 3 . In this embodiment, the heater lead wire passage 25 is formed in an outer peripheral side edge part of the cover 20 . (Fixing Structure of Motor Lead Wires 7 and Heater Lead Wires 8 and the Like) [0049] In order to extend the heater lead wires 8 toward the other side “X2” in the “X” direction and lead out to the outside of the case 3 , at least one of the cover 20 and the case 3 is formed with a lead wire support part for supporting midway portions of the heater lead wires 8 from an outer face side. In this embodiment, the cover 20 is structured with a cover side lead wire support part 26 which supports midway portions of the heater lead wires 8 from an outer face side between an outer face of the case body part 32 of the case 3 and the cover side lead wire support part 26 . More specifically, the heater lead wires 8 are extended in the “X” direction along the side plate part 321 of the case 3 , and the cover side lead wire support part 26 covers the heater lead wires 8 from the other side “Z2” (outer face side) in the “Z” direction and supports the midway portions of the heater lead wires 8 between the side plate part 321 of the case 3 and the cover side lead wire support part 26 . In this embodiment, also in a case that the motor lead wires 7 are to be extended toward the other side “X2” in the “X” direction and led out to the outside of the case 3 , the cover side lead wire support part 26 supports midway portions of the motor lead wires 7 between the outer face of the case body part 32 of the case 3 and the cover side lead wire support part 26 . More specifically, the motor lead wires 7 are extended in the “X” direction along the side plate part 321 of the case 3 , and the cover side lead wire support part 26 covers the motor lead wires 7 from the other side “Z2” (outer face side) in the “Z” direction and supports the midway portions of the motor lead wires 7 between the side plate part 321 and the cover side lead wire support part 26 . Therefore, at the time of assembling of the damper device 1 , portions of the motor lead wires 7 and the heater lead wires 8 extended to the other side “X2” in the “X” direction are fixed by the case 3 and the cover 20 . Therefore, much troublesome labor is not required for fixing portions of the motor lead wires 7 and the heater lead wires 8 extended to the other side “X2” in the “X” direction. [0050] In this embodiment, a portion of the side plate part 321 overlapped with the motor lead wires 7 and the heater lead wires 8 is formed in a recessed part 33 which is recessed to a side of the drive chamber 60 and the cover side lead wire support part 26 is formed in a plate shape. Therefore, when the case 3 is to be fitted to the cover 20 from the other side “X2” in the “X” direction, the cover side lead wire support part 26 is inserted to the recessed part 33 from one side “X1” in the “X” direction and the motor lead wires 7 and the heater lead wires 8 are supported between the bottom part 330 of the recessed part 33 and the cover side lead wire support part 26 . [0051] The cover side lead wire support part 26 is provided with a flat plate part 261 , which covers the motor lead wires 7 and the heater lead wires 8 on the other side “Z2” in the “Z” direction, and protruded plate parts 262 protruded to one side “Z1” in the “Z” direction from both ends in the “Y” direction of the flat plate part 261 . Respective tip ends of the two protruded plate parts 262 are bent toward opposite sides to each other. On the other hand, the recessed part 33 of the case 3 are formed with groove-shaped engaging parts 331 extended in the “X” direction at both ends in the “Y” direction of the bottom part 330 . Therefore, when the cover side lead wire support part 26 is inserted to the recessed part 33 , respective tip ends of two protruded plate parts 262 are fitted to the groove-shaped engaging parts 331 formed in the recessed part 33 of the case 3 . As a result, the cover side lead wire support part 26 is positioned in the “Y” direction and the “Z” direction. Therefore, the position of the cover side lead wire support part 26 is determined surely and thus, the motor lead wires 7 and the heater lead wires 8 can be surely fixed. [0052] In this state, a distance between the side plate part 321 and the cover side lead wire support part 26 (distance between the bottom part 330 of the recessed part 33 and the flat plate part 261 ) is larger than a diameter of the motor lead wire 7 and a diameter of the heater lead wire 8 . Therefore, in a state shown in FIG. 2( b ) , when the cover 20 and the case 3 are to be connected with each other, the motor lead wires 7 and the heater lead wires 8 are not scratched. Accordingly, the cover 20 and the case 3 can be easily connected with each other and insulation coating layers of the motor lead wire 7 and the heater lead wire 8 are hard to be damaged. [0053] Further, in this embodiment, the side plate part 321 of the case 3 is formed with a case side lead wire support part 34 which covers the heater lead wire passage 25 from the other side “Z2” in the “Z” direction and supports midway portions of the motor lead wires 7 from the other side “Z2” (outer face side) in the “Z” direction. [0054] The heater lead wire passage 25 is provided with a first passage which opens toward one side “X1” in the “X” direction, a second passage which is bent from an end part of the first passage on the other side “X2” in the “X” direction to the other side “Y2” in the “Y” direction, and a third passage which is bent from an end part of the second passage on an opposite side to a side of the first passage to the other side “X2” in the “X” direction. The cover side lead wire support part 26 is provided in the third passage. Therefore, after the heater lead wires 8 are pushed to the heater lead wire passage 25 from the other side “Z2” in the “Z” direction, when the heater lead wires 8 are fitted to one side “Z1” in the “Z” direction of the cover side lead wire support part 26 from one side “Y1” in the “Y” direction, the heater lead wires 8 are held by the cover 20 . In this state, even when the heater lead wires 8 are pulled to the other side “X2” in the “X” direction, the pulling force is received by the bent portion of the heater lead wire passage 25 . Therefore, the force is hard to reach to the connecting position 95 of the heater lead wires 8 with the heater 9 and thus connection of the heater lead wires 8 and the heater 9 is hard to be disengaged. (Leading-Around Structure of Motor Lead Wire 7 ) [0055] FIGS. 3( a ), 3( b ) and 3( c ) are explanatory views showing a leading-around structure of motor lead wires 7 in the geared motor 10 and the damper device 1 to which the present invention is applied. FIG. 3( a ) is a plan view showing an inside of the case 3 viewed from one side “X1” in the “X” direction, FIG. 3( b ) is an explanatory view showing a bent structure of the motor lead wires 7 viewed in the “Z” direction, and FIG. 3( c ) is an explanatory view showing a juxtaposed state of the motor lead wires 7 viewed in the “Y” direction. FIGS. 4( a ), 4( b ) and 4( c ) are explanatory views showing an inside of the case 3 in the geared motor 10 and the damper device 1 to which the present invention is applied. FIG. 4( a ) is a perspective view showing an inside of the case 3 , FIG. 4( b ) is a perspective view showing a state that the gear train 65 and the like are detached, an FIG. 4( c ) is an explanatory view showing a state that the motor lead wires 7 are further detached. FIG. 5 is an explanatory view showing a motor lead wire passage 30 which is formed in the case 3 in the geared motor 10 and the damper device 1 to which the present invention is applied. In FIGS. 3( a ), 3( b ) and 3( c ) and FIGS. 4( a ), 4( b ) and 4( c ) , the heater lead wires 8 are not shown. [0056] As shown in FIGS. 3( a ), 3( b ) and 3( c ) , FIGS. 4( a ), 4( b ) and 4( c ) and FIG. 5 , in the geared motor 10 and the damper device 1 in this embodiment, the case 3 is provided with a motor lead wire passage 30 which is extended in the “Y” direction so as to pass a position interposed in the “Z” direction between a turning center axial line of a gear included in the gear train 65 and the case body part 32 . A plurality of the motor lead wires 7 is extended from the other side “Y2” to one side “Y1” in the “Y” direction passing through the motor lead wire passage 30 . In this embodiment, a depth direction of the motor lead wire passage 30 is the “X” direction, and the other side “X2” in the “X” direction of the motor lead wire passage 30 is structured of the bottom plate part 31 . Therefore, the motor lead wire passage 30 is reached to the bottom plate part 31 on the other side “X2” in the “X” direction. [0057] In this embodiment, the case 3 is provided with the motor lead wire passage 30 at a position interposed in the “Z” direction between a turning center axial line “L 68 ” of the drive gear 68 included in the gear train 65 and the side plate part 321 located on the other side “Z2” in the “Z” direction of the case body part 32 . More specifically, the bottom plate part 31 of the case 3 is formed with a support shaft 38 which turnably supports the drive gear 68 and a first plate-shaped projection 351 whose plate thickness direction is the “Z” direction is formed between the support shaft 38 and the side plate part 321 . The first plate-shaped projection 351 is extended in the “Y” direction so as to face the side plate part 321 in the “Z” direction. Therefore, the motor lead wire passage 30 in a slit shape which opens toward the “X” direction is formed between the first plate-shaped projection 351 and the side plate part 321 . Further, in this embodiment, the bottom plate part 31 of the case 3 is also formed with a second plate-shaped projection 352 whose plate thickness direction is the “Z” direction at a position separated from the first plate-shaped projection 351 on the other side “Y2” in the “Y” direction. The second plate-shaped projection 352 is, similarly to the first plate-shaped projection 351 , extended in the “Y” direction so as to face the side plate part 321 in the “Z” direction. Therefore, in this embodiment, the motor lead wire passage 30 in a slit shape which opens toward the “X” direction is formed between the first plate-shaped projection 351 and the side plate part 321 and between the second plate-shaped projection 352 and the side plate part 321 . [0058] Therefore, a plurality of the motor lead wires 7 is extended from the other side “Y2” in the “Y” direction to one side “Y1” passing the motor lead wire passage 30 and, in this state, a part of the plurality of the motor lead wires 7 is overlapped with the drive gear 68 on the other side “X2” in the “X” direction part in the motor lead wire passage 30 when viewed in the “X” direction. [0059] Therefore, in this embodiment, as shown in FIG. 3( b ) and FIG. 4( b ) , a plurality of the motor lead wires 7 is bent in the motor lead wire passage 30 so that at least a part of the motor lead wires 7 is located on a side of the bottom plate part 31 with respect to their connecting positions 79 with the motor 61 and is passed through the motor lead wire passage 30 and extended to one side “Y1” from the other side “Y2” in the “Y” direction. [0060] More specifically, each of a plurality of the motor lead wires 7 is provided with a first portion 71 extended from the connecting position 79 with the motor 61 , a second portion 72 which is bent to the other side “X2” in the “X” direction from an end part of the first portion 71 and extended in the motor lead wire passage 30 , and a third portion 73 which is bent to one side “Y1” in the “Y” direction from an end part on the other side “X2” in the “X” direction of the second portion 72 and is extended in the motor lead wire passage 30 . Further, each of a plurality of the motor lead wires 7 is provided with a fourth portion 74 which is bent from an end part on one side “Y1” in the “Y” direction of the third portion 73 to one side “X1” in the “X” direction and is extended in the motor lead wire passage 30 and a fifth portion 75 which is extended from an end part on one side “X1” in the “X” direction of the fourth portion 74 . Therefore, in each of a plurality of the motor lead wires 7 , at least a part of the third portion 73 is overlapped with the drive gear 68 in the “X” direction when viewed in the “X” direction. This state is maintained by firmness (shape retaining force) of the motor lead wire 7 itself. Further, bent portions in the “X” direction are existed on both sides of the third portion 73 . Therefore, the third portion 73 is hard to be floated to one side “X1” in the “X” direction. [0061] The first portion 71 is extended to the other side “Z2” in the “Z” direction from the connecting position 79 with the motor 61 . Further, the fifth portion 75 is extended to the other side “Z2” in the “Z” direction from the fourth portion 74 and further extended to the other side “X2” in the “X” direction. In this embodiment, the fifth portion 75 is passed through a cut-out part 329 formed in the side plate part 321 of the case 3 and a tip end side of the fifth portion 75 is further extended to the other side “X2” in the “X” direction in a state that the fifth portion 75 is grasped and held from both sides in the “Y” direction by the side plate part 321 of the case 3 . Therefore, even when the fifth portion 75 is pulled to the other side “X2” in the “X” direction, the third portion 73 is hard to be floated to one side “X1” in the “X” direction. Accordingly, even in a case that a plurality of the motor lead wires 7 and the drive gear 68 are overlapped with each other when viewed having looked from “X” direction, the motor lead wires 7 and the drive gear 68 are hard to be contacted with each other. [0062] Further, the cover side projection 28 shown in FIG. 2( b ) is protruded from the cover 20 toward the other side “X2” in the “X” direction at a position displaced to the other side “Y2” in the “Y” direction with respect to the drive gear 68 , and prevents the third portion 73 of the motor lead wire 7 from displacing to one side “X1” in the “X” direction. [0063] In this embodiment, a plurality of the motor lead wires 7 is extended in the “Y” direction in the motor lead wire passage 30 so that they are arranged in parallel or juxtaposed in the “X” direction and, in this state, a plurality of the motor lead wires 7 is structured of the second portion 72 , the third portion 73 , the fourth portion 74 and the like which are bent at plural positions. In order to attain this structure, in this embodiment, a plurality of motor lead wires 7 is structured of a flat cable 70 in which a plurality of motor lead wires 7 are connected with each other in a juxtaposed state in a row. A part in a length direction of the flat cable 70 is separated into a plurality of the motor lead wires 7 and connected with the motor 61 . Also in this case, at least a portion accommodated in the motor lead wire passage 30 is in a state of the flat cable 70 and a thickness direction of the flat cable 70 is the “Z” direction. Therefore, the third portion 73 of a plurality of the motor lead wires 7 is easily accommodated in the motor lead wire passage 30 . [0064] In this embodiment, the third portions 73 of a plurality of the motor lead wires 7 are, as shown in FIG. 3( c ) , juxtaposed in the “X” direction and in a state of the flat cable 70 in an obliquely inclined state in the “Z” direction. Therefore, a region occupied by a plurality of the motor lead wires 7 in the “X” direction is narrow and thus, even in a case that the third portions 73 and the drive gear 68 are overlapped with each other when viewed in the “X” direction, the motor lead wires 7 and the drive gear 68 are hard to be contacted with each other. [0065] In order to attain the structure, in this embodiment, as shown in FIG. 5 , an inner face of the side plate part 321 of the case 3 is formed with a first projection 36 in a plate shape provided with a lead wire receiving part 361 which is obliquely inclined from the “X” direction to the “Z” direction in the motor lead wire passage 30 . The third portions 73 (flat cable 70 ) of a plurality of the motor lead wires 7 are supported in an obliquely inclined state in the “Z” direction along the lead wire receiving part 361 (see FIG. 3( c ) ). Further, an inner face of the side plate part 321 of the case 3 is formed with a second projection 37 in a plate shape provided with a lead wire receiving part 371 which is obliquely inclined from the “X” direction to the “Z” direction in the motor lead wire passage 30 at a position separated from the first projection 36 on one side “Y1” in the “Y” direction. The third portions 73 (flat cable 70 ) of a plurality of the motor lead wires 7 are supported in an obliquely inclined state in the “Z” direction along the lead wire receiving part 371 (see FIG. 3( c ) ). [0000] (Connecting Structure of Motor Lead Wire 7 with Motor 61 ) [0066] In this embodiment, a plurality of the motor lead wires 7 has the same length as each other. The motor lead wire 7 is structured of a core wire having conductivity which is coated by a film made of insulating material and films of adjacent motor lead wires 7 are fusion-bonded to structure the flat cable 70 . A plurality of the motor lead wires 7 is respectively shown by the reference signs 7 A, 7 B, 7 C and 7 D (see FIGS. 3( a ) and 3( b ) ). The motor lead wires 7 A, 7 B, 7 C and 7 D are separated from each other in the first portions 71 extended from the connecting positions 79 with the motor 61 and in the connected portions of the first portions 71 with the second portions 72 . [0067] FIG. 6 is a perspective outward appearance view showing the motor 61 . The motor 61 includes a rotor 62 provided with a permanent magnet on an outer peripheral side of a rotation shaft and a tube shaped stator 63 which surrounds an outer peripheral side of the rotor 62 . The motor 61 is disposed in the drive chamber 60 in a state that a turning center axial line “L 62 ” of the rotor 62 is directed in the “X” direction. The stator 63 includes a pair of outer stator cores 631 and 632 in a bottomed tube shape which are also served as a motor case. The stator 63 is formed with an opening part 633 which is formed by cutting out side faces of the outer stator cores 631 and 632 on the other side “Z2” in the “Z” direction, and a terminal block 64 is provided in the opening part 633 . The terminal block 64 holds a plurality of terminal pins 5 . The number of arranged terminal pins 5 corresponds to the number of the motor lead wires 7 . A plurality of the terminal pins 5 is respectively shown by the reference signs 5 A, 5 B, 5 C and 5 D. The terminal pins 5 A, 5 B, 5 C and 5 D are arranged in a line in this order from one side “Y1” in the “Y” direction toward the other side “Y2”. [0068] The terminal block 64 is extended in the “Y” direction and is provided with a first face 641 which faces the other side “Z2” in the “Z” direction and a second face 642 which faces one side “X1” in the “X” direction. First terminal parts 51 which are one end parts of the terminal pins 5 A, 5 B, 5 C and 5 D are protruded from the first face 641 of the terminal block 64 side by side in a line in the “Y” direction. Further, second terminal parts 52 which are the other end parts of the terminal pins 5 A, 5 B, 5 C and 5 D are protruded from the second face 642 of the terminal block 64 side by side in a line in the “Y” direction. The first terminal parts 51 are portions around which lead-out wires of coils provided in the stator 63 are bound and connected, and the second terminal parts 52 are portions with which core wires of the motor lead wires 7 A, 7 B, 7 C and 7 D are connected by soldering or the like. [0069] As shown in FIG. 3( a ) , the first portions 71 of the motor lead wires 7 ( 7 A, 7 B, 7 C and 7 D) are respectively extended in the “Z” direction in a state that the motor lead wires 7 A, 7 B, 7 C and 7 D are separated from each other and, in addition, these four wires are juxtaposed in this order from one side “Y1” to the other side “Y2” in the “Y” direction. At the connecting positions 79 , the motor lead wire 7 A is connected with the terminal pin 5 A, the motor lead wire 7 B is connected with the terminal pin 5 B, the motor lead wire 7 C is connected with the terminal pin 5 C, and the motor lead wire 7 D is connected with the terminal pin 5 D. [0070] In the second portions 72 , the four motor lead wires 7 A, 7 B, 7 C and 7 D are bent from a state juxtaposed in the “Y” direction to a state juxtaposed in the “X” direction. Further, in the third portions 73 , as described above, the four motor lead wires 7 A, 7 B, 7 C and 7 D are juxtaposed in this order from one side “X1” in the “X” direction to the other side “X2” to structure the flat cable 70 and are extended in the “Y” direction. The fourth portions 74 are bent from the state that the four motor lead wires 7 A, 7 B, 7 C and 7 D are juxtaposed in the “X” direction to a reverse arrangement state to the first portions 71 , in other words, to a state that the motor lead wires 7 D, 7 C, 7 B and 7 A are juxtaposed from one side “Y1” in the “Y” direction toward the other side “Y2” in this order. In the fifth portions 75 , the motor lead wires 7 D, 7 C, 7 B and 7 A are extended in the “Z” direction toward the cut-out part 329 of the case 3 in a state juxtaposed in the “Y” direction. [0071] In this embodiment, work for incorporating the motor lead wires 7 in the motor lead wire passage 30 is, for example, performed as follows. First, before incorporating the motor 61 in the drive chamber 60 , the motor lead wires 7 ( 7 A, 7 B, 7 C and 7 D) and the second terminal parts 52 of the terminal pins 5 ( 5 A, 5 B, 5 C and 5 D) are connected with each other by soldering or the like and, after that, the motor 61 is incorporated in the drive chamber 60 . After that, while successively bending a plurality of the motor lead wires 7 ( 7 A, 7 B, 7 C and 7 D) connected with the terminal pins 5 ( 5 A, 5 B, 5 C and 5 D) so as to structure the first portions 71 , the second portions 72 , the third portions 73 , the fourth portions 74 and the fifth portions 75 , the plurality of the motor lead wires 7 are successively incorporated in the case 3 from the other side “Y2” toward one side “Y1” in the “Y” direction. Principal Effects in this Embodiment [0072] As described above, in the damper device 1 and the geared motor 10 in this embodiment, the case 3 is provided with the motor lead wire passage 30 which is extended in the “Y” direction so as to pass a position interposed in the “Z” direction between the turning center axial line “L 68 ” of the drive gear 68 included in the gear train 65 and the side plate part 321 of the case body part 32 . Therefore, a plurality of the motor lead wires 7 can be extended in the “Y” direction in a state that the positions in the “Z” direction of the motor lead wires 7 are restricted. Further, the plurality of the motor lead wires 7 is extended in the motor lead wire passage 30 from the other side “Y2” to one side “Y1” in the “Y” direction in a state that at least parts of the motor lead wires 7 are bent to a side of the bottom plate part 31 with respect to the connecting positions 79 with the motor 61 , and this state is maintained by firmness (shape retaining force) of the motor lead wire 7 itself. Therefore, the motor lead wires can be led around to a position separated in the “Y” direction from the motor 61 in a state that the motor lead wires 7 and the gear train 65 are not contacted with each other. Accordingly, flexibility for a power feeding position to the motor can be enhanced. [0073] The motor lead wire passage 30 is formed in a slit shape which opens toward one side “X1” in the “X” direction and thus the motor lead wires 7 can be put in the motor lead wire passage 30 from one side “X1” in the “X” direction. Therefore, the motor lead wires 7 are easily led around. [0074] The motor lead wire passage 30 is structured by the bottom plate part 31 on the other side “X2” in the “X” direction. Therefore, parts of the motor lead wires 7 can be bent to a position contacting with the bottom plate part 31 or a position close to the bottom plate part 31 . Therefore, when viewed in the “X” direction, even in a case that a plurality of the motor lead wires 7 and the drive gear 68 are overlapped with each other, the motor lead wires 7 and the drive gear 68 are hard to be contacted with each other. [0075] Parts of a plurality of the motor lead wires 7 are, when viewed in the “X” direction, overlapped with the drive gear 68 on the other side “X2” in the “X” direction in the motor lead wire passage 30 and thus a dimension in the “Z” direction of the geared motor 10 can be reduced. [0076] The third portions 73 of the motor lead wires 7 which are bent toward the other side “X2” in the “X” direction are overlapped with the drive gear 68 in the “X” direction and thus the motor lead wires 7 and the drive gear 68 are hard to be contacted with each other. Further, the third portions 73 of the motor lead wires 7 are prevented from displacing to one side “X1” in the “X” direction by the cover side projection 28 formed in the cover 20 and thus the motor lead wires 7 and the drive gear 68 are hard to be contacted with each other. [0077] The third portions 73 of the motor lead wires 7 are juxtaposed in the “X” direction in the motor lead wire passage 30 and thus a region occupied by the third portions 73 of the motor lead wires 7 is narrow in the “Z” direction. Therefore, even when the third portions 73 of the motor lead wires 7 are disposed at a position near the turning center axial line “L 68 ” of the drive gear 68 , the motor lead wires 7 and the drive gear 68 are hard to be contacted with each other. [0078] The motor 61 is held by the case 3 and the gear train 65 is turnably supported by the bottom plate part 31 of the case 3 . Therefore, when the motor 61 and the gear train 65 are incorporated in the case 3 and then the cover 20 is fitted, assembling work can be easily performed. Further, the cover 20 is integrally structured with the frame 2 and thus, in comparison with a case that the cover 20 is separately structured from the frame 2 , assembling efficiency is improved and the number of components can be reduced. (Modified Embodiment of Leading-Around Structure of Motor Lead Wire 7 ) [0079] In the embodiment described above, a plurality of the motor lead wires 7 ( 7 A, 7 B, 7 C and 7 D) is extended in the “Y” direction in a state juxtaposed in the “X” direction in the motor lead wire passage 30 . However, as shown in FIG. 3( a ) , a plurality of the terminal pins 5 ( 5 A, 5 B, 5 C and 5 D) with which the motor lead wires 7 A, 7 B, 7 C and 7 D are connected is disposed in a line in the “Y” direction and distances of the respective terminal pins from an outlet 30 a of the motor lead wire passage 30 are different from each other. In this case, the outlet 30 a of the motor lead wire passage 30 is the outlet 30 a going toward the other side “Y2” in the “Y” direction from the motor lead wire passage 30 and is provided at a position between an end part on the other side “Y2” in the “Y” direction of the second plate-shaped projection 352 by which the motor lead wire passage 30 is sectioned and the side plate part 321 . [0080] In a modified embodiment, a plurality of motor lead wires 17 (hereinafter, shown by the reference signs 17 A, 17 B, 17 C and 17 D) having lengths different from the above-mentioned embodiment is used. Further, a bending tendency giving process in which the motor lead wires 17 ( 17 A, 17 B, 17 C and 17 D) are previously formed in a shape of a wiring space within the case 3 is performed. Next, the same reference signs are used in the same portions as the embodiment described above and their descriptions are omitted and different portions will be described below by using different reference signs. [0081] FIGS. 7( a ), 7( b ) and 7( c ) are explanatory views schematically showing motor lead wires 17 in a modified embodiment. FIG. 7( a ) shows a state that the motor lead wires 17 are not connected with the motor 61 , FIG. 7( b ) shows a state that the motor lead wires 17 are connected with the motor 61 , and FIG. 7( c ) shows a state that a bending tendency giving process has been performed on the motor lead wires 17 . As shown in FIG. 7( b ) and FIG. 7( c ) , in the modified embodiment, a plurality of motor lead wires 17 A, 17 B, 17 C and 17 D whose lengths are different from each other are connected with the motor 61 . The motor lead wires 17 A, 17 B, 17 C and 17 D are, similarly to the embodiment described above, provided with a portion structuring a flat cable 70 and a portion (separated part 76 ) separated into respective wires. [0082] As shown in FIG. 7( a ) , a plurality of the motor lead wires 17 A, 17 B, 17 C and 17 D are juxtaposed to each other in this order. The separated part 76 is structured so that lengths of the motor lead wires 17 A, 17 B, 17 C and 17 D from a portion connected with the flat cable 70 to their tip ends are set longer in the order of the motor lead wires 17 A, 17 B, 17 C and 17 D. When the separated part 76 is extended toward the motor 61 side, the tip end parts of the motor lead wires 17 A, 17 B, 17 C and 17 D are arranged in this order from one side “Y1” to the other side “Y2” in the “Y” direction. In other words, the tip end parts of the motor lead wires 17 A, 17 B, 17 C and 17 D are disposed in the same order as the terminal pins 5 A, 5 B, 5 C and 5 D on the other side “Z2” in the “Z” direction of the terminal block 64 . [0083] In the modified embodiment, the terminal pins 5 A, 5 B, 5 C and 5 D provided in the motor 61 are sequentially connected with the motor lead wires 17 A, 17 B, 17 C and 17 D in the shorter order of the length and in the nearer order from the outlet 30 a of the motor lead wire passage 30 (in the order of the terminal pins 5 A, 5 B, 5 C and 5 D). As a result, in a state that the motor 61 is incorporated in the case 3 , a plurality of the terminal pins 5 A, 5 B, 5 C and 5 D are connected with the motor lead wires 17 having shorter lengths as they are nearer to the outlet 30 a of the motor lead wire passage 30 . [0084] The motor lead wires 17 A, 17 B, 17 C and 17 D are juxtaposed in this order from the opening part of the motor lead wire passage 30 toward the bottom part (from one side “X1” toward the other side “X2” in the “X” direction) when incorporated in the motor lead wire passage 30 . Therefore, the motor lead wires 17 A, 17 B, 17 C and 17 D are overlapped with each other from one side “X1” to the other side “X2” in the “X” direction in the shorter order of the length at the outlet 30 a of the motor lead wire passage 30 . Accordingly, a plurality of the terminal pins 5 A, 5 B, 5 C and 5 D is sequentially connected with the motor lead wires 17 A, 17 B, 17 C and 17 D in the nearer order from the outlet 30 a of the motor lead wire passage 30 (in the order of the terminal pins 5 A, 5 B, 5 C and 5 D) and in the order of the position located on one side “X1” in the “X” direction. [0085] In the modified embodiment, as described above, the lengths of the motor lead wires 17 A, 17 B, 17 C and 17 D are different from each other and thus, when their end parts are respectively connected with the terminal pins 5 A, 5 B, 5 C and 5 D, the separated part 76 is curved in a shape shown in FIG. 7( b ) as a whole. In other words, the separated part 76 is deformed in a shape bent to one side “Z1” in the “Z” direction as a whole by performing connecting work with the terminal pins 5 A, 5 B, 5 C and 5 D without forcibly bending the separated part 76 to one side “Z1” in the “Z” direction after connection. This shape is a shape corresponding to a shape of the wiring space where the first portions 71 and the second portions 72 are disposed. Therefore, connecting work of the motor lead wires 17 A, 17 B, 17 C and 17 D with the motor 61 is easily performed, and incorporating to the motor lead wire passage 30 and incorporating to the wiring space from the motor lead wire passage 30 toward the drive chamber 60 side are also easily performed. Accordingly, a working hour for wiring work can be shortened. In accordance with an embodiment of the present invention, it may be structured that tip end parts of the motor lead wires 17 A, 17 B, 17 C and 17 D are previously bent and all of four wires are arranged in shapes bent to one side “Z1” in the “Z” direction and, after that, the tip end parts of the motor lead wires 17 A, 17 B, 17 C and 17 D are connected with the terminal pins 5 A, 5 B, 5 C and 5 D. [0086] The motor lead wires 17 A, 17 B, 17 C and 17 D are connected with a plurality of the terminal pins 5 A, 5 B, 5 C and 5 D so that a motor lead wire having a shorter length and located on the opening part side (one side “X1” in the “X” direction) in the motor lead wire passage 30 is connected with the terminal pin nearer to the outlet 30 a of the motor lead wire passage 30 . Therefore, in the first portions 71 and the second portions 72 , it can be avoided that lengths of the motor lead wires 17 A, 17 B, 17 C and 17 D become needlessly longer and that margin of the length is insufficient. For example, excessive slack of the motor lead wire 17 A connected with the terminal pin 5 A which is the nearest to the outlet 30 a can be avoided. Further, the motor lead wire 17 D connected with the terminal pin 5 D on the farthest side from the outlet 30 a can be avoided from being pulled at the time of being incorporated. Further, a longer motor lead wire is led around from the bottom part side of the motor lead wire passage 30 to a far side and a shorter motor lead wire is led around from the opening part side of the motor lead wire passage 30 to a near side and thus the motor lead wires 17 A, 17 B, 17 C and 17 D can be orderly led around. Also from this point, incorporating of the motor lead wires 17 A, 17 B, 17 C and 17 D to the motor lead wire passage 30 and incorporating to the wiring space from the motor lead wire passage 30 toward the drive chamber 60 side are easily performed. Therefore, a working hour for wiring work can be shortened. [0087] As shown in FIG. 7( c ) , in the modified embodiment, after a plurality of the motor lead wires 17 A, 17 B, 17 C and 17 D whose lengths are different from each other is connected with the motor 61 , a bending tendency in a shape bent to the other side “Z2” in the “Z” direction so as to correspond to the shape of the wiring space of the case 3 is given to connecting portions of the fourth portions 74 with the fifth portions 75 of the motor lead wires 17 A, 17 B, 17 C and 17 D in the embodiment described above. For example, the motor lead wires 17 A, 17 B, 17 C and 17 D are wound around a bending tendency giving member 84 in a bar shape as shown by the broken line in FIG. 7( c ) . Next, coating films coated on the core wires of the motor lead wires 17 A, 17 B, 17 C and 17 D are pressed against an outer peripheral face of the bending tendency giving member 84 and the coating films are deformed in a shape corresponding to its outer peripheral face. In this manner, as shown in FIG. 7( c ) , the motor lead wires 17 A, 17 B, 17 C and 17 D are given with a bending tendency in a bent shape. In this case, a jig used for giving a bending tendency is not limited to a bar shape and another shape around which the motor lead wires 17 A, 17 B, 17 C and 17 D are capable of being wound in a bent shape may be adopted. [0088] As described above, when the motor lead wires 17 A, 17 B, 17 C and 17 D are previously given with a bending tendency in a shape of the wiring space in the case 3 , in a case that the bent portions are to be incorporated to the wiring space, bending work in the shape of the wiring space is not required and it is sufficient that the bent portions are fitted to the wiring space. Therefore, workability of wiring work is improved and a working hour can be shortened. [0089] FIGS. 8( a ), 8( b ) and 8( c ) are explanatory views showing an example of a jig which is used for connecting work of the motor lead wires 17 A, 17 B, 17 C and 17 D and the motor 61 . FIG. 8( a ) is a front view, FIG. 8( b ) is a side view, and FIG. 8( c ) shows a working condition. In the modified embodiment, a stator 63 to which a rotor 62 is not assembled is attached to a jig 80 , motor lead wires 17 A, 17 B, 17 C and 17 D are connected with terminal pins 5 A, 5 B, 5 C and 5 D of the stator 63 in advance before assembled into the case 3 . [0090] The jig 80 is provided with a pedestal part 82 which is placed on a workbench 90 , a vertical plate part 81 standing up toward an upper side from the pedestal part 82 , and a stator mounting part 83 in a columnar shape which is protruded from a surface of the vertical plate part 81 . As shown in FIG. 8( c ) , the stator 63 is attached to the stator mounting part 83 in a posture that the terminal block 64 is located on a lower side. In this state, the motor lead wires 17 A, 17 B, 17 C and 17 D are extended to a lower side of the stator mounting part 83 and the separated part 76 is disposed to a lower side of the terminal block 64 . Then, one end parts of the motor lead wires 17 A, 17 B, 17 C and 17 D are lifted one by one to an upper side and connected with the terminal pins 5 A, 5 B, 5 C and 5 D by soldering or the like. [0091] When a bending tendency giving member 84 in a bar shape is provided next to the stator mounting part 83 as shown by the broken line in FIG. 8( c ) , it can be also structured that the motor lead wires 17 A, 17 B, 17 C and 17 D are led around from a lower side of the stator 63 to a side of the bending tendency giving member 84 to perform bending tendency giving work. In this case, it is desirable that a positional relationship between the bending tendency giving member 84 and the stator mounting part 83 is set so as to correspond to a positional relationship between the motor 61 and the bent portions (connected portions of the fourth portions 74 with the fifth portions 75 ) of the motor lead wires 17 A, 17 B, 17 C and 17 D when assembled in the case 3 . [0092] In the modified embodiment, work for incorporating a plurality of the motor lead wires 17 to the motor lead wire passage 30 is, for example, performed as follows. First, motor lead wires 17 A, 17 B, 17 C and 17 D are connected with a stator 63 by using the above-mentioned jig 80 before a rotor 62 is incorporated (first process). Next, portions which structure the fourth portions 74 and the fifth portions 75 of the motor lead wires 17 A, 17 B, 17 C and 17 D are given with a bending tendency in a shape bent by a bending tendency giving member 84 in advance (second process). Next, the rotor 62 and the stator 63 are assembled to structure a motor 61 and then the motor 61 is assembled to a drive chamber 60 of a case 3 (third process). After that, the motor lead wires 17 A, 17 B, 17 C and 17 D which are connected with the terminal pins 5 A, 5 B, 5 C and 5 D of the motor 61 and are given with a bending tendency are incorporated to a wiring space in the case 3 (fourth process). [0093] Required times for wiring work and the like were measured in a case that the motor lead wires 7 A, 7 B, 7 C and 7 D in the embodiment described above are used and a case that the motor lead wires 17 A, 17 B, 17 C and 17 D in the modified embodiment are used. As a result, a total required time of work for incorporating the motor 61 to the drive chamber 60 of the case 3 and work for incorporating the motor lead wires 7 A, 7 B, 7 C and 7 D/ 17 A, 17 B, 17 C and 17 D which are connected with the motor 61 to the wiring space in the case 3 including the motor lead wire passage 30 and leading out from the cut-out part 329 to the outside of the case 3 (in other words, a required time of the third process and the fourth process) was 16.63 seconds in a case of the motor lead wires 7 A, 7 B, 7 C and 7 D in the embodiment described above (lengths are the same and a bending tendency is not given) and was 10.03 seconds in a case of the motor lead wires 17 A, 17 B, 17 C and 17 D in the modified embodiment (lengths are different and a bending tendency is given). Therefore, in the modified embodiment, it was confirmed that a working hour of the wiring work can be shortened. [0094] In the modified embodiment, all of the lengths of a plurality of the motor lead wires 17 A, 17 B, 17 C and 17 D are different from each other but lengths of some of a plurality of the motor lead wires may be different from others. For example, lengths of adjacent two or three motor lead wires may be set the same as each other. Even in this structure, some of a plurality of the terminal pins 5 A, 5 B, 5 C and 5 D can be connected with motor lead wires whose length is shorter as nearer to the outlet 30 a of the motor lead wire passage 30 . Therefore, some of a plurality of the motor lead wires 17 A, 17 B, 17 C and 17 D are avoided from occurring excessive slack and from being pulled due to insufficient margin of the length. Further, the number of the motor lead wires is not limited to four but the present invention may be applied to a structure that motor lead wires having another number are led around. Other Embodiments [0095] Although the present invention has been shown and described with reference to a specific embodiment, various changes and modifications will be apparent to those skilled in the art from the teachings herein. For example, the present invention may be applied to a case that the cover 20 is separately structured from the frame 2 . Further, the damper device 1 in the embodiment described above is for a refrigerator but the present invention is not limited to a damper device used for a refrigerator. REFERENCE SIGNS LIST [0000] 1 damper device 2 frame 3 case 4 baffle 5 , 5 A through 5 D terminal pin 6 baffle drive mechanism 7 , 7 A through 7 D motor lead wire 17 , 17 A through 17 D motor lead wire 8 heater lead wire 9 heater 10 geared motor 20 cover 21 end plate part 22 body part 23 seal plate part 25 heater lead wire passage 26 cover side lead wire support part 28 cover side projection 29 shaft part 30 motor lead wire passage 30 a outlet 31 bottom plate part 32 case body part 33 recessed part 34 case side lead wire support part 36 first projection 37 second projection 38 support shaft 39 hook mechanism 40 opening/closing plate 41 flat plate part 49 elastic member 60 drive chamber 61 motor 62 rotor 63 stator 64 terminal block 65 gear train 66 first gear 67 second gear 68 drive gear 69 sector gear 70 flat cable 71 first portion 72 second portion 73 third portion 74 fourth portion 75 fifth portion 76 separated part 79 connecting position of motor lead wire with motor 80 jig 81 vertical plate part 82 pedestal part 83 stator mounting part 84 bending tendency giving member 95 connecting position of heater lead wire with heater 210 opening part 261 flat plate part 262 protruded plate part 321 through 324 side plate part 329 cut-out part formed in side plate part of case 330 bottom part 331 engaging part 351 first plate-shaped projection 352 second plate-shaped projection 361 lead wire receiving part 371 lead wire receiving part 631 , 632 outer stator core 633 opening part 641 first face 642 second face “L” turning center axial line of baffle “L 68 ” turning center axial line of rotor

The present invention provides a geared motor in which a motor lead can be routed within a case to a position set apart from a motor, and a damper device. In the geared motor that is used in the damper device, the case is provided with a motor lead passage that extends in the Y-direction through a position flanked, in the Z-direction, by a side plate part of a case trunk part and the rotation axis of a driving gear that is included in a gear train. It is therefore possible, with the Z-direction positions of a plurality of motor leads being defined, to extend the motor leads in the Y-direction. The Z-direction dimension of the geared motor can then be reduced because the motor leads will then partially overlap the drive gear in the motor lead passage as viewed from the X-direction.

BACKGROUND AND SUMMARY OF THE INVENTION This invention relates to a circuit and method to reduce the effects of temperature on potentiometer based position sensors. For devices which sense the position of a movable mechanical device or portion thereof, a potentiometer or variable resistor is often used. The wiper arm of the potentiometer is connected to the movable portion of the mechanical device such that motion of the device results in motion of the wiper arm. The changed resistance results in a signal which corresponds to the amount of motion of the mechanical device. When this type of position sensor is employed in a high temperature environment, the potentiometer may exhibit a drop in resistance with increasing temperature. This change in resistance could be interpreted by the electronic monitor as a change in the position of the mechanical device. This false reading may result in an improper or untimely command being processed by the control circuitry. The subject invention is designed to minimize the thermal effects on the position sensor. The scheme involves measuring the time needed for a comparator to change states given different input signals. The change-of-state time period under one input condition is divided by the time period under a second input condition. The quotient of the time periods is a function of resistance values which translates to a measure of the actuator's position. Both sides of the potentiometer are subject to the effects of temperature. These effects are cancelled out by taking the quotient of the time periods. DESCRIPTION OF THE DRAWINGS Other objects, features and advantages of the present invention will become more fully apparent from the following detailed description of the preferred embodiment, the appended claims and the accompanying drawings in which FIG. 1 is a schematic of an electronic circuit used to practice the preferred embodiment of the temperature compensation method described herein. FIG. 2 is an alternative embodiment of the circuit shown in FIG. 1. FIG. 3 is a voltage vs. time graph of selected points in the circuit shown in FIG. 1. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENT Referring to FIG. 1, the subject invention includes comparator means with the non-inverting input V + set to a value of 1/3 of the supply voltage Vcc. The inverting input V - to the comparator is presented with a fluctuating signal, the time of the fluctuation being determined by the presence or absence of a control signal at nodes A and B. The fluctuation of the signal to the inverting input V - to comparator U 1 is also affected by the position of the wiper arm on a potentiometer R. The potentiometer's wiper arm is connected physically to the movable portion of a mechanical device. The motion of the movable portion of the mechanical device will then reflect itself in changing resistance values of R 1 and R 2 . The preferred embodiment shown in the FIG. 1 sets up comparator U 1 with its output being sent to a controller such as a microprocessor. The controller is able to store the output curve for comparator U 1 . The non-inverting input V + to comparator U 1 is set as a reference level. In this case the reference level is 1/3 Vcc. The inverting input V - to U 1 receives its voltage level from a network which is made up of resistor R 4 connected between the supply voltage Vcc and the inverting input V - , resistor R 3 , and a discharging loop connected to node D. Resistor R 3 is connected between the inverting input V - to U 1 and the discharging loop at node D. The discharging loop is made up of a switching transistor Q 1 which receives its input command from the controller at node A, a potentiometer R and a capacitor C which receives its stimulus from the controller at node B. The capacitor C AC couples node B and node D. The potentiometer R is made up of two resistance portions R 1 and R 2 which vary in value depending on the location of the wiper arm. Also to be accounted for is the wiper arm contact resistance which is represented as a resistance R w connected between the wiper arm contact point to potentiometer R and node D. The wiper arm is electronically connected to node D and physically connected to the movable portion of the mechanical device. The R 1 portion of potentiometer R is connected between ground and the wiper arm; the R 2 portion of potentiometer R is connected between the wiper arm and the collector of switching transistor Q 1 . The emitter of switching transistor Q 1 is connected to ground and the ihput signal from the controller to switching transistor Q 1 is made at node A to the base of the transistor. The operation of the circuit is as follows: Referring to FIG. 2 and FIG. 3, with the voltage V A at node A in a low state, transistor Q 1 is turned off and the discharging loop includes R 1 but not R 2 . This is the case since transistor Q 1 blocks current flow from R 2 through Q 1 to ground. The voltage V A at node A is a controlled event, preferably from a control microprocessor. The voltage V B at node B is also a controlled event. With transistor Q 1 turned off, and with the voltage V B at node B turned off, voltage V D at node D and at the inverting input V - to U 1 is held low resulting in the output voltage V 0 of comparator U 1 being on. The above described events all have occurred between time zero and time t 1 on the voltage vs. time curves in FIG. 3. At time t 1 , the controller raises the voltage V B at node B to Vcc. Since capacitor C AC couples nodes B to node D the voltage V D will spike to Vcc (with a slight overshoot). The duration of this control pulse is unimportant so long as it resets to zero before time t 4 . Immediately after time t 1 , the voltage at the inverting input V - crosses the fixed threshold voltage (1/3 Vcc) at the non-inverting input V + , turning comparator U 1 off. The voltage at node D will peak at approximately Vcc and then begin to discharge through resistor R 1 . This is the R 1 discharge path. At time t 2 , the voltage at the inverting input V - discharges down to the threshold voltage of 1/3 Vcc turning comparator U 1 on. The discharge time t 2 -t 1 is stored by the controller as the representation of the R 1 path discharge time. At time t 3 , the controller changes the voltage V A at node A to a high state turning transistor Q 1 on and bringing R 2 in parallel with R 1 in the discharge path. This is the "R 1 , R 2 path". At time t 4 , the controller pulses node B again, repeating the discharge events but this time utilizing the "R 1 , R 2 path" for discharge. The output V o of comparator U 1 turns off at time t 4 when the threshold voltage is exceeded at the V - input to U 1 and the comparator U 1 turns back on at time t 5 when V - drops below the threshold voltage 1/3 Vcc. The time period (t 5 -t 4 ) is stored by the controller and represents the "R 1 , R 2 path" discharge time. The charging times are: (t 2 -t 1 )=-R 1 C×ln (1/3) and (t 5 -t 4 )=-[R 1 R 2 /(R 1 +R 2 )] C ln (1/3). The quotient of the times is: (t 2 -t 1 )/(t 5 -t 4 )=R 2 /(R 1 +R 2 ) and is a measure of the position of the movable mechanical device. The quotient of the discharge times can be inverted and still be a measure of the position of the movable portion of the mechanical device. The quotient or its inverse must be used consistently throughout the analysis. It can be seen that the error due to temperature effects is a function of the position of the wiper arm on potentiometer R. The error will be reflected unevenly in the numerator and denominator of the quotient. The temperature affects both arms of the potentiometer and the arm with the largest resistance will have a largest error. This affects the quotient since only one arm (R 2 ) is affected in the numerator of the quotient (or the denominator if the inverse quotient is used). Although this effect should be noted, it is not always significant; the temperature compensation provided by the circuit and method described above often provides enough compensation. For some situations temperature effects can be dealt with more effectively by using the alternative embodiment shown in FIG. 2. In this embodiment an additional transistor Q 2 is placed into the circuit between R 1 and ground, the collector of Q 2 connected to R 1 and the emitter of Q 2 to ground. The base of Q 2 is available for control signals from the controller at A'. The method is then changed to measure the (t 2 -t 1 ) time by keeping Q 1 off and turning Q 2 on, placing R 1 in the discharge path. The discharge time (t 2 -t 1 ) is measured as before. Then Q 2 is turned off and Q 1 is turned on placing only R 2 in the discharge path. The discharge time (t 5 -t 4 ) is measured as before. The result generates (t 2 -t 1 ) as a function of R 1 and (t 5 -t 4 ) as a function of R 2 . The resulting quotient (R 2 /R 1 ) will still contain an error which is a function of the resistance of each arm of potentiometer R; the larger the resistance, the larger the error. While the present invention has been disclosed in connection with the preferred embodiment thereof, it should be understood that there may be other embodiments which fall within the spirit and scope of the invention and that the invention is susceptible to modification, variation and change without departing from the proper scope or fair meaning of the following claims.

A circuit and method to compensate for thermal effects on potentionmeter based sensors by presenting two input signals at different times to a comparator; the first utilizing one leg of the potentiometer and the second utilizing both legs. The comparator output will change states after receiving each input signal. The time periods between the receipt of the input signal and the outputs' change of state are stored. The quotient of the time periods cancels out most of the thermal effects on the potentiometer and is still a representative of the original potentiometer signal.

BACKGROUND OF THE INVENTION The present invention relates to an electromechanical watch including a time base, a frequency divider, a stepping motor driving the indicating members of the watch, a maintaining circuit for the motor controlled by the frequency divider, and correcting means manually operated by way of a control member, allowing action on some of the indicating members. Such watches are known. The above mentioned correcting means allow, for instance, action on only the date indicating member, or else only on the hour indicating member, without modifying the position of any other indicating member. This provides the user with the possibility of correcting the indication of the date after months having thirty days, or of correcting the indication of the time when travelling and passing from one time zone to another. In the last case, the correcting means are generally arranged in such a way as to allow the advance of the hour indicator step by step at the rate of one hour for each step. The drawback of such watches lies in the fact that the indicating members, of the minute and of the second, for instance, do not occupy a position which is absolutely exact after the correction of the hour indicator which is particularly disagreeable in the case of electronic watches, the precision of which is such that one desires not to lose either the exact minute, or the exact second. The purpose of the present invention is to overcome this drawback by furnishing overtaking means which causes the indicating members of the watch to resume, after their correction, a position corresponding exactly to the actual time, the watch being stopped during the correction period. SUMMARY OF THE INVENTION The electro-mechanical watch according to the invention is characterized by the fact that it includes a pulse counter, and a switching device controlled by the control member which is arranged to interrupt the pulsing of the stepping motor and to direct to the counter the pulses coming from the frequency divider during the correction period. A logic circuit controls a maintaining circuit for the motor and is arranged in such a way that, when the correcting operation is ended, the motor is then driven at a speed higher than its normal speed a number of steps equal to the number of pulses the counter has received, which restores the normal state of the indicating members, and compensates for the loss of steps occuring during the correction period. BRIEF DESCRIPTION OF THE DRAWINGS The drawing shows, by way of example, one embodiment of the invention. FIG. 1 is a diagrammatic plan view of a portion of an electro-mechanical watch showing its control mechanism which allows the conventional setting of the watch and the correction of the hour indicator only. FIG. 2 is a sectional view of a detail, along line II--II of FIG. 1. FIG. 3 is a diagram of the electronic circuit of this watch, and FIG. 4 is a diagram showing the state of different points of the circuit of FIG. 3. DESCRIPTION OF THE PREFERRED EMBODIMENT The watch represented in the drawing includes an electric stepping motor 1, only the shaft 1' of which, has been represented in FIG. 1. This motor drives, by the intermediary of two wheels 2 and 3, an hour cannon wheel 4 in FIG. 2. The hour hand, designated by 5, is not mounted directly on the hour cannon wheel 4, as in conventional constructions, but on a cannon 6a of a supplementary cannon wheel 6 which is itself engaged on a cannon 4a of the hour cannon wheel 4. The cannon 4a is axially split, so as to be resilient, and has, at its end, an outer shoulder 4b cooperating with an inner shoulder of the cannon 6a for maintaining the wheel 6 in place. A resilient washer 7 is interposed between the wheels 4 and 6 and is provided with protrusions 7a directed towards the wheel 4 and with protrusions 7b directed towards the wheel 6. The protrusions 7a engage holes of the wheel 4 while the protrusions 7b engage holes of the wheel 6. This resilient washer 7 thus ensures the angular connection between the wheels 4 and 6. The frusto-conical shape of the protrusions 7a have the effect that, when the torque to be transmitted between these two wheels goes beyond a predetermined value, the protrusions 7a leave the holes of the wheel 4 and a slipping movement occurs. The wheel 4 is meshed with a pinion, designated by 8, of a dial-train the wheel of which, designated by 9, is meshed with a pinion 10a of the cannon-pinion, designated by 10, carrying the minute hand 11 in FIG. 2. The watch as represented includes a control stem 12 in FIG. 1 which is able to occupy two axial positions, pushed and pulled. Stem 12 when pushed, as represented in the drawing, selects the function it is desired to be effected and when pulled, effects the selected function. As represented in the example, the stem 12 in the pulled position, is able to effect two functions, i.e. on one hand the conventional setting of the time indicating members, such as the hour hand 5 and minute hand 11, or even the second hand, and on the other hand the driving or setting of the hours-hand 5 only. To effect these two functions, the stem 12 carries a pinion 13 meshing with an intermediate pinion 14 frictionally mounted on a rocking lever 15, and coaxial with an axis 16 of the articulation of this lever on the frame of the movement. Hence, when the stem 12 occupies its pushed position, represented in FIG. 1, a rotation of stem 12 in one direction or the other causes the lever 15 to rotate in one direction or the other, as indicated by arrows 17 and 18, which produces the selection of the desired function. When the stem 12 is pulled, it brings a beak or finger 19a of a setting lever, 19, which is articulated at 20 on the frame of the movement, to come in front of or behind a beak or finger 15a of the lever 15, thus maintaining lever 15 in the selected position. The lever 15 carries two intermediate pinions 21 and 22 meshing with the intermediate pinion 14. Depending on the selection which has been made when stem 12 is in the pushed position, one or the other of the two following situations occurs: (a) The intermediate pinion 21 is meshing with the wheel 9 of the dial-train which is itself meshing with teeth 10a of the cannon pinion 10. Hence, the rotation of the stem 12, in one direction or the other, allows one to effect the conventional setting of the time indicating members, such as the hour hand 5 and the minute hand 11, or even the second hand; and (b) The intermediate pinion 22 is meshing with the lower portion of a double intermediate pinion 23 the upper portion of which is meshing with the supplementary hour cannon wheel 6 carrying the hour hand 5. The rotation of the stem 12 then allows one to act only on the hour hand, without moving the minute hand or the second hand. Because the protrusions 7a of the resilient washer 7 ensure the driving of the wheel 4 by the wheel 6, the driving of the hour hand only is effected step by step, at the rate of one hour for each step. Such mechanisms are known per se, one of them appearing, for instance, in the Swiss Pat. No. 571,736. The lever 15 is moreover acted on by two return blade springs 24 and 25 cooperating with the shaft, 26 of the intermediate pinion 21, and which springs tend to maintain lever 15 in its rest position as represented in FIG. 1. When the lever 15 occupies one or the other of its two working positions, one or the other of the springs 24 and 25 is resiliently deformed. When spring 24 is deformed, it comes into contact with a pin 27, and when spring 25 is deformed it comes into contact with a pin 28. The elastic spring 24 and pin 27, elastic spring 25 and pin 28 constitute respectively, switches 24-27 and 25-28. The electronic circuit represented in FIG. 3 includes a time base 29 constituted for example, by a piezo-electric resonator, a frequency divider 30, a maintaining circuit 31, the motor 1, a counter 32, the two switches 24-27 and 25-28, two OR gates 33 and 34, three AND gates 35, 36 and 37, and three inverters 38, 39 and 40. The elements contained in the block indicated by 41 constitute a switching circuit, while those contained in the block indicated by 42 constitute a logic circuit. In normal operation, the frequency of the output signal of the time base 29 is divided by the divider 30, the frequency of the output 30A of this divider being applied to the input 35a of the AND gate 35. Since the input 35b of the AND gate 35 is at the logic state 1, because the switches 24-27 and 25-28 are open, the output signal of the divider 30 is gated to the output of AND gate 35. This signal is applied to the maintaining circuit 31 through the OR gate 34. The maintaining circuit 31 produces from this signal shaped pulses which maintain or drive the motor 1. The other elements of the circuit then do not play any role. One may determine, indeed, that the AND gate 36 is locked or disabled and will not pass any pulses as a result of the switch 24-27 being open, which causes the input 33a and output 33a of the OR gate 33 to be at the logic state 1, which causes, due to the inverter 38, the input 36a of the AND gate 36 to be at the logic state 0. The input 32a of the counter 32 does not therefore receive any pulsed signal. When the count of counter 32 is zero, its output 32A is at the logic state 1. Due to the inverter 40, the input 37a of the AND gate 37 is at the logic state 0, so that this gate is also disabled. If one effects a correction of the hour hand, as a result in a change in the time zone, for instance, one closes the switch 24-27, as indicated previously. The input 33a of the OR gate 33 is then at the logic state 0. If, at this moment, the divider 30 emits a pulse, the input 33b of the OR gate 33 is at the logic state 1 up to the end of this pulse. The output 33A then passes to the logic state 0, which brings the input 36a of the AND gate 36 to the logic state 1 due to the inverter 38. In this way, each time a pulse is emitted from the frequency divider 30, it is passed to the output 36A of the AND gate 36, wherefrom it increments the counter 32. The AND gate 37 remains disabled as a result of its input 37b being at the logic state 0. The AND gate 35 is also disabled as a result of its input 35b being at the logic state 0. When, at the end of the correction period, the switch 24-27 is re-opened, the input 33a of the OR gate 33 passes to the logic state 1. Consequently, the AND gate 35 is enabled to allow passage of the pulses coming from the frequency divider 30. The output 32A of the counter 32 is no longer at the logic state 1 since the content of this counter is not zero. The input 37a of the AND gate 37 is at the logic state 1. The input 37b is at the logic state 1 since the switch 24-27 is open. As long as the divider 30 does not emit a pulse, the input 37c of the AND gate 37 is also at the logic state 1, due to the inverter 39. The switch 25-28 being also open, the input 37d of the AND gate 37 is at the logic state 1. In this way, pulses coming from an intermediary output 30B of the frequency divider 30 are passed to the output 37A of the AND gate 37. On one hand, the pulses from output 37A are subtracted from the content or count of the counter 32, and on the other hand, they are applied to the driving circuit 31 of the motor through the OR gate 34. The pulses from the output 37A are of higher frequency than those coming from the output 30A of the divider 30 which cause the motor 1 to advance at a speed higher than its normal speed, which compensates for the lack of advancement of the indicating members of the watch occuring during the correction period. If, during this compensation period, the divider 30 delivers a pulse through its output 30A, the input 37c of the AND gate 37 passes to the logic state 0, in such a way that a pulse coming simultaneously from the output 30B of the divider 30 is not passed to the output 37A of the AND gate 37. Consequently, the pulse from output 30A does not act on the counter 32, but it is passed to the output 35A of the AND gate 35, so that it acts on the maintaining circuit 31. When the counter 32 again has a count equal to zero, its output 32A passes to the logic state 1, which brings the input 37a of the AND gate 37 to the logic state 0, due to the inverter 40. The watch then resumes its normal running or operation. If the user wants to effect the conventional setting of the watch, he closes, during the manipulation, the switch 25-28, which is the usual switch for electro-mechanical watches. Therefore, during conventional setting, the motor is stopped, its pulsing being interrupted, and the divider is brought back to zero, so that, when the stem is pushed into its rest position, the motor starts after a time equal to one step, which moves the second hand, at the desired time. The closing of the switch 25-28 has the effect of resetting to zero the divider 30 and the counter 32. The motor 1 is no longer driven. The input 36a of the AND gate 36 is also at the logic state 0. Consequently, the input 32a of the counter 32 does not receive any signal. During this period, there is no action, either on the motor, or on the counter. The closing of the switch 25-28 has a further effect of resetting to zero the divider 30 and the counter 32 in such a way that, when the user, after having ended the setting of the watch, pushes the control stem 12, the watch starts again normally. It is to be noted that, so as to ensure the stability of the positions of the rotor of the motor when the switch 25-28 is closed, the setting lever 19 carries a resilient whip 43 (FIG. 1) which cooperates, when the stem 12 is in its pulled position, with one or the other of two flat surfaces 44 of the shaft 1 of the motor. FIG. 4 shows the state of different points of the circuit of FIG. 3: Line a corresponds to the output 30A of the divider 30; line b corresponds to the output 30B of the divider; line c corresponds to the output of the switch 25; line d corresponds to the output of the switch 24; line e corresponds to the output 35A of the AND gate 35; line f indicates the pulses applied to the motor 1; line g corresponds to the input 36a of the AND gate 36; line h corresponds to the input 32a of the counter 32; line i corresponds to the output 32A of the counter 32; line j corresponds to the input 37c of the AND gate 37; line k corresponds to the output 37A of the AND gate 37.

An electro-mechanical watch for advancing the uncorrected watch indicating members to their correct time indicating positions, after a correction period during which one of the indicating members is adjusted and the advancement of the uncorrected members is interrupted. In regular operation, a pulsed signal at a normal frequency from a watch frequency divider drives a stepping motor to advance the indicating members. At the beginning of the correction period, the pulsed signal is switched from the stepping motor to a pulse counter which counts the pulses occurring during the correction period. After the correction period, the pulsed signal is switched back to the stepping motor and a number of signal pulses at a higher than normal frequency, equal to the count of the pulse counter are applied to the stepping motor to advance rapidly the uncorrected indicating members to their correct time indicating positions.

This invention relates to multigrade lubricating oils for use in lubricating internal combustion engines, that contain basestocks with low levels of saturated hydrocarbons, and specifically to such oils which also comprise a multifunctional viscosity modifier. Multigrade lubricating oils typically are identified by designations such as SAE 10W-30, 5W-30 etc. The first number in the multigrade designation is associated with a maximum low temperature (e.g.,-20° C.) viscosity requirement for that multigrade oil as measured typically by a cold cranking simulator (CCS) under high shear rates (ASTM D5293, which is a revision of ASTM D2602), while the second number in the multigrade designation is associated with a high temperature viscosity requirement usually measured in terms of the kinematic viscosity (kV) at 100° C. (ASTM D445). Thus, each particular multigrade oil must simultaneously meet both strict low and high temperature viscosity requirements, set e.g. by SAE specifications such as SAE J300, in order to qualify for a given multigrade oil designation. The high temperature viscosity requirement is intended to prevent the oil from thinning out too much during engine operation which can lead to excessive wear and oil consumption. The maximum low temperature viscosity requirement is intended to facilitate engine starting in cold weather and to ensure pumpability, i.e., the cold oil should readily flow to the oil pump, otherwise the engine can be damaged due to insufficient lubrication. The viscosity characteristic of a basestock on which a lubricating oil is based is typically expressed by the neutral number of the oil (e.g., S150N) with a higher neutral number being associated with a higher viscosity at a given temperature. Blending basestocks is one way of modifying the viscosity properties of the resulting lubricating oil. Unfortunately, merely blending basestocks of different viscosity characteristics may not enable the formulator to meet the low and high temperature viscosity requirements of some multigrade oils. The formulator's primary tool for achieving this goal is an additive conventionally referred to as a viscosity modifier (VM) or viscosity index improver (V.I. improver). A monofunctional VM is conventionally an oil-soluble long chain polymer. A multifunctional VM (or alternately MFVM) is an oil soluble polymer which has been chemically modified e.g., functionalized and derivatized, to impart dispersancy as well as viscosity modification. The basestocks which are typically used in lubricating oils may be synthetic or natural oils. Mineral oils contain various amounts of saturated hydrocarbons, such as straight or branched chain paraffins and naphthenes, and unsaturated hydrocarbons particularly aromatic hydrocarbons. Lubricating oils have traditionally used basestocks containing high levels of saturated hydrocarbon--also referred to as high saturate basestocks--since aromatic hydrocarbons give rise to difficulties in formulating for adequate performance in internal combustion engines. This has been known for some time, being discussed, for example, in "Lubricants for Fluid Film and Hertzian Contact Conditions", T. I. Fowle, Proc. Instn. Mech. Engrs. 1967-8, Vol 182, Pt 3A, pages 568-576, especially pages 568/9 and 571/2. More recently, "Chemistry and Technology of Lubricants", edited by R. M. Mortier and S. T. Orszulik, Blackie Academic and Professional, 1992, in chapter 1, "Base oils from Petroleum" R. J. Prince, pages 1-31, discusses the instability of aromatic components to oxidation which is still perceived as a problem. "Compositional Analysis of Re-refined and Non-Conventional Lubricant Base Oils: Correlations to Sequence VE and IIIE Gasoline Engine Tests", Stipanovic et aL, SAE Technical Paper Series, 941978, Oct. 17-20 1994 provides a statistical analysis in those engine tests which indicates a strong negative impact of various aromatic hydrocarbon types. Among other consequences it is generally accepted that there is a tendency for unsaturated components and particularly aromatic components of basestocks to contribute to the formation of baked-on deposits in engines, generally referred to as "varnish". As discussed in the literature identified above, special and expensive finishing treatments are required to remove aromatics from basestocks and so increase the level of saturates. Increasingly there is a need for lubricating oils for internal combustion engines which are capable of utilising basestocks with low levels of saturates. In order to meet stringent engine performance requirements and specifically to give adequate varnish inhibition to those oils with conventional types of additive formulations it has proved necessary to use very high treat levels of dispersants and/or to use specific detergent systems. This is economically undesirable and also give rise to further problem within the oil formulation, as those high levels of additives can bring their own problems of oxidation stability, compatibility and engine performance debits. This invention relates to multigrade lubricating oils which utilise low saturate basestocks and provide adequate varnish performance without requiring high levels of dispersant and/or detergent additives. Thus, in one aspect the invention provides a multigrade lubricating oil for an internal combustion engine which comprises: a. a basestock of lubricating oil viscosity having less than 75 mass % of saturated hydrocarbons; b. less than 3 mass % of ashless dispersant derived from a polymer of number average molecular weight (Mn) of not greater than 5000; and c. viscosity modifier to give the desired viscometrics, which comprises at least one multifunctional viscosity modifier. DETAILED DESCRIPTION A. Basestock As indicated above conventional lubricating oils are prepared using basestocks which have relatively high levels of saturates and thus low levels of unsaturated and specifically aromatic hydrocarbons. Mineral basestocks are typically subjected to hydrogen treatments such as hydrocracking or hydroisomerisation in order to give greater paraffinic content and lower aromatic content. The basestock used in the lubricating oil of the invention does not require such treatments and may use lower grade basestocks previously regarded as unsuitable for such applications. Such basestocks for use in the invention are typically mineral oils which have not been subjected to severe treatments to raise the saturates level, but the invention may employ any of the available synthetic or natural oils, re-refined oils and mixtures of such oils, provided the overall saturates level of the basestock or basestock mixture is less than 75 mass %, preferably less than 70 mass %, and may even use basestocks of less than 65 mass % saturates. Such basestocks may contain at least 20%, preferably at least 30 mass % of aromatic compounds and may even contain in excess of 35 mass % of aromatic compounds. Additives used in formulating lubricating oils often contain diluent oil; this diluent oil introduced with additives is not included within the term "basestock" as that term is used herein, which is confined to the oil used to dilute the additives to form the finished oil. The lubricating oil basestock conveniently has a viscosity of from 2.5 to 12 mm 2 /s, and preferably from 2.5 to 9 mm 2 /s, at 100° C. Examples of commercially available basestocks of low saturates content which may be employed in the invention are ESN 600 (typically 69.9 mass % saturates; 30.1 mass % aromatics) available from Esso Petroleum Co. Ltd., Agip 450 (typically 64.7 mass % saturates; 35.3 mass % aromatics) available from Agip Petroli and BP 500ME (typically 61.9 mass % saturates; 38.1 mass % aromatics) available from B.P. pic. Such low saturate basestocks may be used alone or in combination with other basestocks, which may also have low saturates content or have relatively higher saturate content, provided that the saturate content of the combined basestock as that term is used herein is less than 75 mass % of the total basestock. B. Ashless Dispersant The ashless dispersant comprises an oil soluble polymeric hydrocarbon backbone having functional groups that are capable of associating with particles to be dispersed. Typically, the dispersants comprise amine, alcohol, amide, or ester polar moieties attached to the polymer backbone often via a bridging group. The ashless dispersant may be, for example, selected from oil soluble salts, esters, amino-esters, amides, imides, and oxazolines of long chain hydrocarbon substituted mono and dicarboxylic acids or their anhydrides- thiocarboxylate derivatives of long chain hydrocarbons; long chain aliphatic hydrocarbons having a polyamine attached directly thereto, and Mannich condensation products formed by condensing a long chain substituted phenol with formaldehyde and polyalkylene polyamine. The oil soluble polymeric hydrocarbon backbone is typically an olefin polymer, especially polymers comprising a major molar amount (i.e. greater .ia 50 mole %) of a C 2 to C 18 olefin (e.g., ethylene, propylene, butylene, isobutylene, pentene, octene-1, styrene), and typically a C 2 to C 5 olefin. The oil soluble pentene, octene-1, styrene), and typically a C 2 to C 5 olefin. The oil soluble polymeric hydrocarbon backbone may be a homopolymer (e.g. polypropylene or polyisobutylene) or a copolymer of two or more of such olefins (e.g. copolymers of ethylene and an alphaolefin such as propylene and butylene or copolymers of two different alpha-olefins). Other copolymers include those in which a minor molar amount of the copolymer monomers, e.g., 1 to 10 mole %, is a C 3 to C 22 non-conjugated diolefin (e.g., a copolymer of isobutylene and butadiene, or a copolymer of ethylene, propylene and 1,4-hexadiene or 5-ethylidene-2-norbornene). One preferred class of olefin polymers is polybutenes and specifically polyisobutenes (PIB) or poly-n-butenes, such as may be prepared by polymerization of a C 4 refinery stream. Another preferred class of olefin polymers is ethylene alpha-olefin (EAO) copolymers or alpha-olefin homo- and copolymers having in each case a high degree (e.g.>30%) of terminal vinylidene unsaturation. That is, the polymer has the structure: P-HCR=CH 2 wherein P is the polymer chain and R is a C 1 -C 18 alkyl group, typically methyl or ethyl. Preferably the polymers have at least 50% of the polymer chains with terminal vinylidene unsaturation. EAO copolymers of this type preferably contain 1 to 50 mass % ethylene, and more preferably 5 to 45 mass % ethylene. Such polymers may contain more than one alpha-olefin and may contain one or more C 3 to C 22 diolefins. Also usable are mixtures of EAO's of low ethylene content with EAO's of high ethylene content. The EAO's may also be mixed or blended with PIB's of various Mn's or components derived from these may be mixed or blended. Atactic propylene oligomer typically having Mn of from 700 to 500 may also be used, as described in EP-A490454. Suitable olefin polymers and copolymers, such as polyisobutenes, may be prepared by cationic polymerization of hydrocarbon feedstreams, usually C 3 -C 5 , in the presence of a strong Lewis acid catalyst and a reaction promoter, usually an organoaluminum such as HCI or ethylaluminum dichloride. Tubular or stirred reactors may be used. Such polymerizations and catalysts are described, e.g., in U.S. Pat. Nos. 4,935,576 and 4,952,739. Fixed bed catalyst systems may also be used as in U.S. Pat. No. 4,982,045 and UK-A 2,001,662. Most commonly, polyisobutylene polymers are derived from Raffinate I refinery feedstreams. Conventional Ziegier-Natta polymerization may also be employed to provide olefin polymers suitable for use to prepare dispersants and other additives. The preferred EAO polymers may be prepared by polymerizing the appropriate monomers in the presence of a catalyst system comprising at least one metallocene (e.g. a cyclopentadienyl-transition metal compound) and preferably an activator, e.g. an alumoxane compound. The metallocenes may be formed with one, two, or more cyclopentadienyl groups, which are substituted or unsubstituted. The metallocene may also contain a further displaceable ligand, preferably displaced by a cocatalyst--a leaving group--that is usually selected from a wide variety of hydrocarbyl groups and halogens. Optionally there is a bridge between the cyclopentadienyl groups and/or leaving group and/or transition metal, which may comprise one or more of a carbon, germanium, silicon, phosphorus or nitrogen atom-containing radical. The transition metal may be a Group IV, V or VI transition metal. Such polymerizations and catalysts are described, for example, in U.S. Pat, Nos. 4,871,705, 4,937,299, 5,017,714, 5,120,867, 4,665,208, 5,153,157, 5,198,401, 5,241,025, 5,057,475, 5,096,867, 5,055,438, 5,227,440, 5,064,802; EP-A-129368, 520732, 277003, 277004, 420436; WO91/04257, 93/08221, 93/08199 and 94/13715. The oil soluble polymeric hydrocarbon backbone of the ashless dispersant, as that term is used herein, has a number average molecular weight (Mn) of not greater than 5,000. The Mn of the backbone is preferably within the range of 500 to 5,000, more preferably 700 to 5,000 where the use of the backbone is to prepare a component having the primary function of dispersancy. Hetero polymers such as polyepoxides are also usable to prepare components. Both relatively low molecular weight (Mn 500 to 1500) and relatively high molecular weight (Mn 1500 to 5,000) polymers are useful to make dispersants. Particularly useful olefin polymers for use in dispersants have Mn within the range of from 1500 to 3000. The Mn for such polymers can be determined by several known techniques. A convenient method for such determination is by gel permeation chromatography (GPC) which additionally provides molecular weight distribution information, see W. W. Yau, J. J. Kirkland and D. D. Bly, "Modern Size Exclusion Liquid Chromatography", John Wiley and Sons, New York, 1979. The oil soluble polymeric hydrocarbon backbone may be functionalized to incorporate a functional group into the backbone of the polymer, or as pendant groups from the polymer backbone. The functional group typically will be polar and contain one or more hetero atoms such as P, 0, S, N, halogen, or boron. It can be attached to a saturated hydrocarbon part of the oil soluble polymeric hydrocarbon backbone via substitution reactions or to an olefinic portion via addition or cycloaddition reactions. Alternatively, the functional group can be incorporated into the polymer by oxidation or cleavage of a small portion of the end of the polymer (e.g., as in ozonolysis). Useful functionalization reactions include. halogenation of the polymer at an olefinic bond and subsequent reaction of the halogenated polymer with an ethylenically unsaturated functional compound. reaction of the polymer with an unsaturated functional compound by the "ene" reaction absent halogenation (an example of the former functionalization is maleation where the polymer is reacted with maleic acid or anhydride); reaction of the polymer with at least one phenol group (this permits derivatization in a Mannich Base-type condensation), reaction of the polymer at a point of unsaturation with carbon monoxide using a Koch-type reaction to introduce a carbonyl group in an iso or neo position, reaction of the polymer with the functionalizing compound by free radical addition using a free radical catalyst, reaction with a thiocarboxylic acid derivative; and reaction of the polymer by air oxidation methods, epoxidation, chioroamination, or ozonolysis. The functionalized oil soluble polymeric hydrocarbon backbone is then further derivatized with a nucleophilic amine, amino-alcohol, or mixture thereof to form oil soluble salts, amides, imides, amino-esters, and oxazolines. Useful amine compounds include those described herein after in more detail in relation to the MFVM. Preferred amines are aliphatic saturated amines. Non-limiting examples of suitable amine compounds include. 1,2-diaminoethane; 1,3-diaminopropane; 1,4-diaminobutane; 1,6-diaminohexane; polyethylene amines such as diethylene triamine; triethylene tetramine; tetraethylene pentamine; and polypropyleneamines such as 1,2-propylene diamine; and di-(1,2-propylene)triamine. Useful amines also include polyoxyalkylene polyamines and the polyamido and related amido-amines as disclosed in U.S. Pat, Nos. 4,857,217, 4,956,107, 4,963,275 and 5229022. Also usable is tris(hydroxymethyl)amino methane (THAM) as described in U.S. Pat. Nos. 4,102,798, 4,113,639 and 4,116,876; and GB-A-989409. Dendrimers, star-like amines, and comb-structure amines may also be used. Similarly, one may use the condensed amines of U.S. Pat. No. 5,053,152. The functionalized polymer of this invention is reacted with the amine compound according to conventional techniques as in EP-A-208560 and U.S. Pat. No. 5,229,022 using any of a broad range of reaction ratios as described therein. A preferred group of nitrogen containing ashless dispersants includes those derived from polyisobutylene substituted with succinic anhydride groups and reacted with polyethylene amines (e.g. tetraethylene pentamine, pentaethylene, polyoxypropylene diamine) aminoalcohols such as trismethylolaminomethane and optionally additional reactants such as alcohols and reactive metals e.g. pentaerythritol, and combinations thereof). Also useful as nitrogen containing ashless dispersants are dispersants wherein a polyamine is attached directly to the long chain aliphatic hydrocarbon as shown in U.S. Pat, Nos. 3,275,554 and 3,565,804 where a halogen group on a halogenated hydrocarbon is displaced with various alkylene polyamines. Another class of nitrogen-containing ashless dispersants comprises Mannich base condensation products. Such Mannich condensation products may include a long chain, high molecular weight hydrocarbon (e.g., Mn of 1,500 or greater) on the benzene group or may be reacted with a compound containing such a hydrocarbon, for example, polyalkenyl succinic anhydride as shown in U.S. Pat. No. 3,442,808. Examples of dispersants prepared from polymers prepared from metallocene catalysts and then functionalized, derivatized, or functionalized and derivatized are described in U.S. Pat. Nos. 5,266,223, 5,128,056, 5,200,103, 5,225,092, 5,151,204 and 5,334,775; WO-A-94/13709 and 94/19436; and EP-A440506, 513211 and 513157. The functionalizations, derivatizations, and post-treatments described in the following patents may also be adapted to functionalize and/or derivative the preferred polymers described above: U.S. Pat. Nos. 3,275,554, 3,565,804, 3,442,808, 3,442,808, 3,087,936 and 3254025. C. Viscosity Modifiers The multifunctional viscosity modifier may be one or more of: polymethacrylates derivatised with nitrogen containing monomers such as vinylpyridine, N-vinylpyrrolidinone, or N,N'-dimethylaminoethyl methacrylate; ethylene-propylene copolymers directly amine derivatised, hydrogenated star polymers reacted with a carboxylic acid derivative and then reacted with an amine; hydrogenated styrenebutadiene-ethylene oxide block copolymers; and ethylene alphaolefin copolymers solution or melt grafted with ethylenically unsaturated a dicarboxylic acid derivative and then reacted with an amine. Typically multifunctional viscosity modifiers are derived from a polymer having a number average molecular weight (Mn) of greater than 7000, as distinct from ashless dispersants, as defined above. In a preferred aspect the multifunctional viscosity modifier comprises a derivatized ethylene-alpha olefin copolymer comprising an adduct of (i) a copolymer having a number average molecular weight of from 20,000 to 100,000, functionalized with mono- or dicarboxylic acid material; and (ii) at least one amine, and in a particularly preferred aspect the ethylene-alpha olefin copolymer comprises either a) from 30 to 60 weight percent monomer units derived from ethylene and from 70 to 40 weight percent monomer units derived from alpha-olefin, or b) from 60 to 80 weight percent monomer units derived from ethylene and from 40 to 20 weight percent monomer units derived from alpha olefin. A highly preferred class of multifunctional viscosity modifiers which may be used in the invention comprise a mixture of derivatised ethylene-alpha olefin copolymers A and B, both comprising an adduct of (i) a copolymer having a number average molecular weight of from 20,000 to 100,000, functionalized with mono- or dicarboxylic acid material; and (ii) at least one amine, and wherein: the ethylene-alpha olefin copolymer of derivatized copolymer A comprises from 30 to 60 weight percent monomer units derived from ethylene and from 70 to 40 weight percent monomer units derived from alpha-olefin; and the ethylene-alpha olefin copolymer of derivatized copolymer B comprises from 60 to 80 weight percent monomer units derived from ethylene and from 40 to 20 weight percent monomer units derived from alpha olefin, with the proviso that the respective weight percents of ethylene derived monomer units present in said derivatized copolymers A and B differ by at least 5 weight percent. The multifunctional viscosity modifiers used in the present invention may be prepared by known techniques. The preferred mixture of derivatized ethylene-alpha-olefin copolymers may be prepared by functionalising and derivatising ethylene alphaolefin copolymers such as described in EP-A-616616 and WO-A-94/1 3763. Ethylene Alpha-olefin Copolymers The ethylene-alpha olefin copolymers comprise monomer units derived from ethylene and alpha-olefins which are typically C 3 to C 28 , preferably C 3 to C 18 , most preferably C 3 to C 8 alpha olefins. While not essential, such polymers preferably have a degree of crystallinity of less than 25 wt. percent as determined by x-ray and differential scanning calorimetry. Copolymers of ethylene and propylene are most preferred. Representative examples of other suitable alpha-olefins include 1-butene, 1-pentene, 1-hexene, 1-heptene, 1-octene, 1-nonene, 1-decene, etc; also branched chain alpha-olefins, such as 4 methyl-1-pentene, 4-methyl-1-hexene, 5 methyl pentene-1, 4.4 dimethyl-1-pentene, and 6 methylheptene-1 and mixtures thereof. Ter- and tetra- copolymers are included within the scope of "copolymers". Ethylene alpha-olefin copolymers used in the invention preferably have a number average molecular weight (Mn) of from 25,000 to 80,000 and most preferably from 25,000 to about 50,000. Suitable polymers will typically have a narrow molecular weight distribution (MWD), as determined by the ratio of weight average molecular weight (Mw) to number average molecular weight (Mn). Polymers having a Mw/Mn of less than 10, preferably less than 7, and more preferably 4 or less are most desirable. As used herein (Mn) and (Mw) may be measured by well known techniques such as vapor phase osmometry (VPO), membrane osmometry and gel permeation chromatography (GPC). The synthesis of polymers having a suitable molecular weight and narrow MWD may be obtained by techniques known in the art including choice of synthesis conditions and post synthesis treatment such as extrusion at elevated temperature, high shear mastication under elevated temperatures in the presence of peroxides or air. thermal degradation, and fractional precipitation from solution. The copolymers employed to make the component blends of the present invention are differentiated primarily by their ethylene content. Derivatised copolymer A is derived from a low ethylene monomer unit content copolymer and derivatised copolymer B is derived from a high ethylene monomer unit content copolymer. More specifically, the low ethylene content copolymer will comprise preferably from 40 to 50 and most preferably from 42 to 46 (e.g., 44) weight percent monomer units derived from ethylene; and preferably from 60 to 50, and most preferably from 58 to 54 (e.g., 56) weight percent monomer units derived from alpha-olefin. The high ethylene content copolymer will comprise preferably from 65 to 75, and most preferably from 68 to 73 (e.g., 70) weight percent monomer units derived from ethylene; and preferably from 35 to 25, and most preferably from 32 to 27 (e.g., 30) weight percent monomer units derived from alpha-olefin. The above ethylene contents are subject to the proviso that the ethylene content of the high and low ethylene copolymers must differ by at least 5, preferably at least and most preferably at least 15 weight percent. For ease of discussion, derivatised copolymers derived from the low ethylene content copolymer, as described above, are referred to herein as Component A, and derivatised copolymers derived from the high ethylene content copolymer, as described above, are referred to herein as Component B. Many such ethylene alpha olefin copolymers are available as items of commerce and their composition and methods for producing them are well known in the art. Representative examples include: MDV-90-9 manufactured by Exxon Chemical Company, an ethylene-propylene copolymer containing 70 weight percent ethylene, which is further characterized by a Mooney viscosity, ML, 1+4 @ 125° C. of 18; and VISTALON 457 manufactured by. Exxon Chemical Company, a 44 weight percent ethylene, ethylene-propylene copolymer which is further characterized by a Mooney viscosity, ML 1+4 @ 125° C. of 28. As indicated above, the MFVM used in present invention comprises a blend of Components A and B. Such blends will comprise typically weight ratios (referred to herein as "blend ratios") of A: B of from 2.3:1 to 0.18: 1, preferably from 1.2:1 to 0.25: 1, and most preferably from 0.8:1 to 0.33:1. Such blend ratios are also applicable to unfunctionalized high and low ethylene content polymer blends in preparation for functionalization. To prepare the MFVM used in the present invention, the high and low ethylene alpha-olefin copolymers are first functionalized and then derivatized. Functionalized Polymers By functionalized, it is meant that the polymer is chemically modified to have at least one functional group present within its structure, which functional group is capable of undergoing further chemical reaction (e.g., derivatization) with other materials. The preferred functionalization reaction is accomplished by reaction of the polymer with a compound containing the desired functional group by free radical addition using a free radical catalyst. More specifically, polymer functionalized with mono- or dicarboxylic acid material, i.e., acid, anhydride, salt or acid ester suitable for use in this invention, typically includes the reaction product of the polymer with a monounsaturated carboxylic reactant comprising at least one of (i) monounsaturated C 4 to C 10 dicarboxylic acids (preferably wherein (a) the carboxyl groups are vicinyl, i.e., located on adjacent carbon atoms and (b) at least one, more preferably both, of said adjacent carbon atoms are part of said monounsaturation). (ii) derivatives of (i) such as anhydrides or C 1 to C 5 alcohol derived mono- or diesters of (i); (iii) monounsaturated C 3 to C 10 monocarboxylic acids wherein the carbon-carbon double bond is conjugated allylic to the carboxyl group, i.e., of the structure --C═C--CO--; and (iv) derivatives of (iii) such as C 1 to C 5 alcohol derived monoesters of (iii). Suitable unsaturated acid materials thereof which are useful functional compounds, include acrylic acid, crotonic acid, methacrylic acid, maleic acid, maleic anhydride, fumaric acid, itaconic acid, itaconic anhydride, citraconic acid, citraconic anhydride, mesaconic acid, glutaconic acid, choromaleic acid, aconitic acid, crotonic acid, methylcrotonic acid, sorbic acid, 3-hexenoic acid, 10-decenoic acid, 2-pentenel,3,5-tricarboxylic acid, cinnamic acid, and lower alkyl (e.g., C 1 to C 4 alkyl) acid esters of the foregoing, e.g., methyl maleate, ethyl fumarate, methyl fumarate, etc. Particularly preferred are the unsaturated dicarboxylic acids and their derivatives, especially maleic acid, fumaric acid and maleic anhydride. The two functionalised copolymers described above can be prepared in several ways. The functional groups can be grafted onto each of the copolymers separately and then the functionalized copolymers can then be mechanically blended at the above described blend ratios. In the preferred method for practicing the invention, the two copolymers are simultaneously functionalized and blended at the same time by feeding into an extruder, masticator or reactor. The extrusion process is continuous, while the masticator process is a batch process. Both take place in a polymer melt, i.e., the polymer is melted in the high temperature, high shear conditions of this equipment. The functionalization takes place substantially in absence of a solvent. The reactor process is a process similar to the masticator batch process but the polymer is functionalized once it is dissolved in a solvent such as mineral oil. The extruder and masticator processes can provide efficient peroxide and or thermo oxidative induced molecular weight reduction of the copolymers, should a lower molecular weight be desired than that of the copolymer that is available. It will be understood that blends of the high and low ethylene content polymers will create a bimodal distribution of ethylene content not achievable by making a single polymer having a single average ethylene content. Free-radical induced grafting can take place in a polymer melt in a extruder or masticator, or when using a conventional batch reactor with the polymer dissolved in a solvent, preferably in a mineral lubricating oil. The free-radical grafting is preferably carried out using free radical initiators such as peroxides, hydroperoxides, and azo compounds and preferably those which have a boiling point greater than about 100° C. and which decompose thermally within the grafting temperature range to provide said free radicals. The initiator is generally used at a level of between about 0.005 percent and about 1 percent, based on the total weight of the polymer. The ethylenically unsaturated carboxylic acid material, preferably maleic anhydride, will be generally used in an amount ranging from 0.01 percent to 10 percent, preferably 0.1 to 2.0 percent, based on weight of copolymer. The aforesaid carboxylic acid material and free radical initiator are generally used in a weight percent ratio range of 1.0:1 to 30:1, preferably 3.0:1 to 6:1. When the copolymer grafting takes place in a solvent in a reactor, the initiator grafting is preferably carried out in an inert atmosphere, such as that obtained by nitrogen blanketing. While the grafting can be carried out in the presence of air, the yield of the desired graft polymer is generally thereby decreased as compared to grafting under an inert atmosphere substantially free of oxygen. The grafting time will usually range from 0.1 to 12 hours, preferably from 0.5 to 6 hours, more preferably 0.5 to 3 hours. In the grafting process, usually the copolymer solution is first heated to grafting temperature and thereafter said unsaturated carboxylic acid material and initiator are added with agitation, although they could have been added prior to heating. When the reaction is complete, the excess acid material can be eliminated by an inert gas purge, e.g., nitrogen sparging. The grafting is preferably carried out in a mineral lubricating oil which need not be removed after the grafting step but can be used as the solvent in the subsequent reaction of the graft polymer with the amine material and as a solvent for the end product to form the lubricating additive concentrate. The oil having attached, grafted carboxyl groups, when reacted with the amine material will also be converted to the corresponding derivatives but such derivatives are of little use to improvement in performance. A description for functionalizing in a masticator can be found in U.S. Pat. No. 4,735,736, and a description for functionalizing the copolymers, dissolved in a solvent such as mineral oil, in a reactor can be found in U.S. Pat. No. 4,517,104, the disclosures of which are herein incorporated by reference. In contrast, reactions carried out in the polymer melt, particularly in an extruder, are characterized by maximized reaction rates and minimized reactor volumes (due to the absence of a diluent solvent), by absence of side reactions with the solvent and by minimized residence times (due to the absence of dissolution and recovery steps before and after the reaction, respectively). Methods for extruder grafting are disclosed in commonly assigned U.S. Pat. No. 5,290,461, the disclosure of which is herein incorporated by reference. In order to prevent or minimize the crosslinking or gellation of the grafted copolymer, particularly when it is subsequently aminated with amines having more than one reactive primary or secondary nitrogens, an optional acid functionalized low molecular weight hydrocarbyl component can be added to the functionalized polymers to moderate molecular weight growth of the derivatized polymer. Such materials are referred to herein as "Growth Regulators". Suitable Growth Regulators include. hydrocarbyl substituted succinic anhydride or acid having 12 to 49 carbons, preferably 16 to 49 carbons in said hydrocarbyl group, long chain monocarboxylic acid of the formula RCOOH where R is a hydrocarbyl group of 50 to 400 carbons and long chain hydrocarbyl substituted succinic anhydride or acid having 50 to 400 carbons in said hydrocarbyl group. Primarily because of its ready availability and low cost, the hydrocarbyl portion, e.g., alkenyl groups, of the carboxylic acid or anhydride is preferably derived from a polymer of a C 2 to C 5 monoolefin, said polymer generally having a molecular weight of about 140 to 6500, e.g., 700 to about 5000, most preferably 700 to 3000 molecular weight. Particularly preferred is polyisobutylene of 950 molecular weight. Derivatized Polymers A derivatized polymer is one which has been chemically modified to perform one or more functions in a significantly improved way relative to the unfunctionalized polymer and or the functionalized polymer. The primary new function sought to be imparted to the functionalized polymers of the present invention is dispersancy in lubricating oil compositions. Thus, the derivatized polymers used in the invention are the reaction products of the above recited functionalized polymers with amines. Of the various amines useful in the practice of this invention, one amine type has two or more primary amine groups, wherein the primary amine groups may be unreacted, or wherein one of the amine groups may already be reacted. Particularly preferred amine compounds include alkylene polyamines, polyoxyalkylene polyamines, preferably wherein the alkylene groups are straight or branched chains containing from 2 to 7, and more preferably 2 to 4 carbon atoms. Examples of the alkylene polyamines include methylene amines, ethylene amines, butylene amines, propylene amines, pentylene amines, hexylene amines, heptylene amines, octylene amines, other polymethylene amines, the cyclic and higher homologs of these amines such as the piperazines, the amino-alkyl-substituted piperazines, etc. These amines include, for example, ethylene diamine, diethylene triamine, triethylene tetramine, propylene diamine, di(heptamethylene)triamine, tripropylene tetramine, tetraethylene pentamine, trimethylene diamine, pentaethylene hexamine, di(trimethylene)triamine, 2-heptyl-3-(2-aminopropyl)imidazoline, 4-methylimidazoline, 1,3-bis-(2-aminoethyl)imidazoline, pyrimidine, 1-(2-aminopropyl)-piperazine, 1,4-bis-(2-aminoethyl)piperazine, N,N-dimethyaminopropyl amine, N,N-dioctylethyl amine, N-octyl-N'-methylethylene diamine, 2-methyl-1-(2-aminobutyl) piperazine, etc. The ethylene amines which are particularly useful are described, for example, in the Encyclopaedia of Chemical Technology under the heading of "Ethylene Amines" (Kirk and Othmer), Volume 5, pgs. 898-905. Interscience Publishers, New York (1950). The polyoxyalkylene polyamines are preferably polyoxyalkylene diamines and polyoxyalkylene triamines, and may typically have average molecular weights ranging from 200 to 4000 and preferably from 400 to 2000. The preferred polyoxyalkylene polyamines include the polyoxyethylene and polyoxypropylene diamines and the polyoxypropylene triamines having average molecular weights ranging from 200 to 2000. The polyoxyalkylene polyamines are commercially available and may be obtained, for example, from the Jefferson Chemical Company, Inc. under the trade name "Jeffamines D-230, D-400, D-1000, D-2000, T-403", etc. Primary amines are more preferred because of the stability of the imide products formed. Most preferred are primary amines, RNH 2 , in which the R group contains functionalities that it is desired to have in the final product. Although such products contain two functionalities, the imide functionality formed by reaction of the primary amine is relatively inert and serves as a stable linkage between the functionality in the R group and the polymer backbone. In this invention it is desired that the R group of the primary amine RNH 2 contain tertiary amine functionality. Examples of useful primary amines, RNH 2 , in which the R group contains tertiary amine functionality include: N,N-dimethylethylenediamine, N,N-diethylethylenediamine, N, N-dimethyl-1,3-propanediamine, N,N-diethyl-1,3-propanediamine, 4-aminomorpholine, 4-(aminomethyl)pyridine, 4-(2-aminoethyl)morpholine and 4-(3-aminopropyl)morpholine. Preferred reactive compounds for reaction with grafted maleic anhydride in the practice of this invention are 4-(3-aminopropyl)morpholine and 1-(2-aminoethyl)- piperazine. Still other amines useful in the practice of this invention include amino-aromatic polyamine compounds such as N-arylphenylenediamines. Particularly preferred N-arylphenylenediamines are the N-phenylphenylenediamines, for example, N-phenyl-1,4-phenyienediamine, N-phenyl-1,3-phenylenediamine, N-phenyl-1,2-phenylenediamine, N-naphthyl-phenylenediamine, N-phenyl-naphthalenediamine and N'-aminopropyl-N-phenylphenylene- diamine. Other useful amines include aminothiazoles such as aminothiazole, aminobenzothiazole, aminobenzothiadiazole and aminoalkylthiazole, aminopyrroles, phenothiazines and phenothiazine derivatives, particularly 10-aminopropyl-phenothiazine, amino-3-propylaminophenothiazine, N-amino-propyl-2-naphthylamine and N-aminopropyidiphenylamine. Mixtures of amines, particularly mixtures of two or more of the above compounds, may be used. As indicated above, functionalization can be conducted separately on the high and low ethylene content polymers or the high and low ethylene content polymers can be blended at the aforedescribed blend ratios and then functionalized. If the latter option is employed, derivatization is conducted on the blend. If separate functionalization is employed, one has the additional options of derivatizing separately and blending the final derivatized products or blending the separately functionalized copolymers and derivatizing the blend simultaneously. The functionalized ethylene alpha-olefin copolymers can be derivatized with amine in the melt or in solution. Melt derivatizations can in turn be conducted in an extruder or masticator, when conditions are substantially the same as the functionalization step. A stripping step can take place prior to amination to remove the unwanted by-products of the graft step which can lead to undesirable by-products as a consequence of the amination. When the amination takes place in a reactor, the functionalized polymer is dissolved in solution (e.g., in oil) at an amount of typically from 5 to 30, preferably 10 to 20, wt. percent polymer, based on the solution weight. Accordingly, the functionalized polymer is preheated at a temperature of from about 100° C. to 250° C., preferably from 170° to 230° C., said amine and optional growth regulator added and temperatures maintained for from 1 to 10 hours, usually 2 to 6 hours. It has been found that many of these multifunctional viscosity modifiers which contain unreacted primary or secondary amine, can undergo an increase in molecular weight which is manifested by product gellation or viscosity growth of the resultant concentrates in oil. For this reason it has been found useful to post-treat or cap these products with an acid such as a C 12 to C 16 hydrocarbyl substituted dicarboxylic acid or anhydride to stabilize the molecular weight. The lubricating oils of the invention typically contain a minor amount, e.g. 0.001 up to 50 mass percent, preferably 0.005 to 25 mass percent, based on the weight of the lubricating oil, of the derivatized copolymers as MFVM. The viscosity modifier system used in the invention will be used in an amount to give the required viscosity characteristics. When used in lubricating oils for automotive or diesel crankcase lubrication the MFVM is present at concentrations usually within the range of from 0.01 to 10 mass percent, e.g., 0.1 to 6.0 mass percent, preferably 0.25 to 3.0 mass percent (measured as polymer), of the total composition. A single multifunctional viscosity modifier may be used alone, or it may be used in combination with additional conventional viscosity modifiers, either monofunctional or multifunctional. Additional additives are typically incorporated into the compositions of the present invention. Examples of such additives are ashless dispersants, metal or ash containing detergents, antioxidants, anti-wear agents, friction modifiers, rust inhibitors, anti-foaming agents, demulsifiers, and pour point depressants. D. Detergent Metal-containing or ash-forming detergents function both as detergents to reduce or remove deposits and as acid neutralizers or rust inhibitors, thereby reducing wear and corrosion and extending engine life. Detergents generally comprise a polar head with a long hydrophobic tail, with the polar head comprising a metal salt of an acidic organic compound. The salts may contain a substantially stoichiometric amount of the metal in which case they are usually described as normal or neutral salts, and would typically have a total base number or TBN (as may be measured by ASTM D2896) of from 0 to 80. It is possible to include large amounts of a metal base by reacting an excess of a metal compound such as an oxide or hydroxide with an acidic gas such as carbon dioxide. The resulting overbased detergent comprises neutralised detergent as the outer layer of a metal base (e.g. carbonate) micelle. Such overbased detergents may have a TBN of 150 or greater, and typically of from 250 to 450 or more. Detergents that may be used include oil-soluble neutral and overbased sulfonates, phenates, sulfurized phenates, thiophosphonates, salicylates, and naphthenates and other oil-soluble carboxylates of a metal, particularly the alkali or alkaline earth metals, e.g., sodium, potassium, lithium, calcium, and magnesium. The most commonly used metals are calcium and magnesium, which may both be present in detergents used in a lubricant, and mixtures of calcium and/or magnesium with sodium. Particularly convenient metal detergents are neutral and overbased calcium sulfonates having TBN of from 20 to 450 TBN, and neutral and overbased calcium phenates and sulfurized phenates having TBN of from 50 to 450. Sulfonates may be prepared from sulfonic acids which are typically obtained by the sulfonation of alkyl substituted aromatic hydrocarbons such as those obtained from the fractionation of petroleum or by the alkylation of aromatic hydrocarbons. Examples included those obtained by alkylating benzene, toluene, xylene, naphthalene, diphenyl or their halogen derivatives such as chlorobenzene, chlorotoluene and chloronaphthalene. The alkylation may be carried out in the presence of a catalyst with alkylating agents having from about 3 to more than 70 carbon atoms. The alkaryl sulfonates usually contain from about 9 to about 80 or more carbon atoms, preferably from about 16 to about 60 carbon atoms per alkyl substituted aromatic moiety. The oil soluble sulfonates or alkaryl sulfonic acids may be neutralized with oxides, hydroxides, alkoxides, carbonates, carboxylate, sulfides, hydrosulfides, nitrates, borates and ethers of the metal. The amount of metal compound is chosen having regard to the desired TBN of the final product but typically ranges from about 100 to 220 mass % (preferably at least 125 mass %) of that stoichiometrically required. Metal salts of phenols and sulfurised phenols are prepared by reaction with an appropriate metal compound such as an oxide or hydroxide and neutral or overbased products may be obtained by methods well known in the art. Sulfurised phenols may be prepared by reacting a phenol with sulfur or a sulfur containing compound such as hydrogen sulfide, sulfur monohalide or sulfur dihalide, to form products which are generally mixtures of compounds in which 2 or more phenols are bridged by sulfur containing bridges. E. Antiwear and Antioxidant Agent Dihydrocarbyl dithiophosphate metal salts are frequently used as anti-wear and antioxidant agents. The metal may be an alkali or alkaline earth metal, or aluminum, lead, tin, molybdenum, manganese, nickel or copper. The zinc salts are most commonly used in lubricating oil in amounts of 0.1 to 10, preferably 0.2 to 2 mass %, based upon the total weight of the lubricating oil composition. They may be prepared in accordance with known techniques by first forming a dihydrocarbyl dithiophosphoric acid (DDPA), usually by reaction of one or more alcohol or a phenol with P 2 S 5 and then neutralizing the formed DDPA with a zinc compound. For example, a dithiophosphoric acid may be made by reacting mixtures of primary and secondary alcohols. Alternatively, multiple dithiophosphoric acids can be prepared where the hydrocarbyl groups on one are entirely secondary in character and the hydrocarbyl groups on the others are entirely primary in character. To make the zinc salt any basic or neutral zinc compound could be used but the oxides, hydroxides and carbonates are most generally employed. Commercial additives frequently contain an excess of zinc due to use of an excess of the basic zinc compound in the neutralization reaction. The preferred zinc dihydrocarbyl dithiophosphates are oil soluble salts of dihydrocarbyl dithiophosphoric acids and may be represented by the following formula. ##STR1## wherein R and R' may be the same or different hydrocarbyl radicals containing from 1 to 18, preferably 2 to 12, carbon atoms and including radicals such as alkyl, alkenyl, aryl, arylalkyl, alkaryl and cycloaliphatic radicals. Particularly preferred as R and R' groups are alkyl groups of 2 to 8 carbon atoms. Thus, the radicals may, for example, be ethyl, n-propyl, i-propyl, n-butyl, i-butyl, sec-butyl, amyl, n-hexyl, i-hexyl, n-octyl, decyl, dodecyl, octadecyl, 2-ethylhexyl, phenyl, butylphenyl, cyclohexyl, methylcyclopentyl, propenyl, butenyl. In order to obtain oil solubility, the total number of carbon atoms (i.e. R and R') in the dithiophosphoric acid will generally be about 5 or greater. The zinc dihydrocarbyl dithiophosphate can therefore comprise zinc dialkyl dithiophosphates. Conveniently at least 50 (mole) % of the alcohols used to introduce hydrocarbyl groups into the dithiophosphoric acids are secondary alcohols. Oxidation inhibitors or antioxidants reduce the tendency of mineral oils to deteriorate in service which deterioration can be evidenced by the products of oxidation such as sludge and varnish-like deposits on the metal surfaces and by viscosity growth. Such oxidation inhibitors include hindered phenols, alkaline earth metal salts of alkylphenolthioesters having preferably C 5 to C 12 alkyl side chains, calcium nonylphenol sulfide, ashless oil soluble phenates and sulfurized phenates, phosphosulfurized or sulfurized hydrocarbons, phosphorous esters, metal thiocarbamates, oil soluble copper compounds as described in U.S. Pat. No. 4,867,890, and molybdenum containing compounds. Typical oil soluble aromatic amines having at least two aromatic groups attached directly to one amine nitrogen contain from 6 to 16 carbon atoms. The amines may contain more than two aromatic groups. Compounds having a total of at least three aromatic groups in which two aromatic groups are linked by a covalent bond or by an atom or group (e.g., an oxygen or sulfur atom, or a --CO--, --SO 2 -- or alkylene group) and two are directly attached to one amine nitrogen also considered aromatic amines. The aromatic rings are typically substituted by one or more substituents selected from alkyl, cycloalkyl, alkoxy, aryloxy, acyl, acylamino, hydroxy, and nitro groups. OTHER ADDITIVES Friction modifiers may be included to improve fuel economy. Oil-soluble alkoxylated mono- and diamines are well known to improve boundary layer lubrication. The amines may be used as such or in the form of an adduct or reaction product with a boron compound such as a boric oxide, boron halide, metaborate, boric acid or a mono-, di- or trialkyl borate. Other friction modifiers include esters formed by reacting carboxylic acids and anhydrides with alkanols. Other conventional friction modifiers generally consist of a polar terminal group (e.g. carboxyl or hydroxyl) covalently bonded to an oleophillic hydrocarbon chain. Esters of carboxylic acids and anhydrides with alkanols are described in U.S. Pat. No. 4,702,850. Examples of other conventional friction modifiers are described by M. Belzer in the "Journal of Tribology" (1992), Vol. 11 4, pp. 675-682 and M. Belzer and S. Jahanmir in "Lubrication Science" (1988), Vol. 1, pp. 3-26. Rust inhibitors selected from the group consisting of nonionic polyoxyalkylene polyols and esters thereof, polyoxyalkylene phenols, and anionic alkyl sulfonic acids may be used. Copper and lead bearing corrosion inhibitors may be used, but are typically not required with the formulation of the present invention. Typically such compounds are the thiadiazoie polysuifides containing from 5 to 50 carbon atoms, their derivatives and polymers thereof. Derivatives of 1,3,4 thiadiazoies such as those described in U.S. Pat. Nos. 2,719,125; 2,719,126, and 3,087,932, are typical. Other similar materials are described in U.S. Pat. Nos. 3,821,236; 3,904,537; 4,097,387; 4,107,059; 4,136,043. 4,188,299. and 4,193,882. Other additives are the thio and polythio sulfenamides of thiadiazoies such as those described in UK. Patent Specification No. 1,560,830. Benzotriazoies derivatives also fall within this class of additives. When these compounds are included in the lubricating composition, they are preferably present in an amount not exceeding 0.2 mass % active ingredient. A small amount of a demulsifying component may be used. A preferred demulsifying component is described in EP 330,522. It is obtained by reacting an alkylene oxide with an adduct obtained by reacting a bis-epoxide with a polyhydric alcohol. The demulsifier should be used at a level not exceeding 0.1 mass % active ingredient. A treat rate of 0.001 to 0.05 mass % active ingredient is convenient. Pour point depressants, otherwise known as lube oil flow improvers, lower the minimum temperature at which the fluid will flow or can be poured. Such additives are well known. Typical of those additives which improve the low temperature fluidity of the fluid are C 8 to C 18 dialkyl fumarate/vinyl acetate copolymers and polyalkylmethacrylates. Foam control can be provided by many compounds including an antifoamant of the polysiloxane type, for example, silicone oil or polydimethyl siloxane. Lubricating compositions may also contain elastomer comparability aids for elastomeric seals such as Viton or fluorocarbon seals and nitrile seals. Carboxylic acids and unsaturated hydrocarbons have been used for such a purpose. Some of the above-mentioned additives can provide a multiplicity of effects; thus for example, a single additive may act as a dispersant-oxidation inhibitor. This approach is well known and does not require further elaboration. When lubricating compositions contain one or more of the above-mentioned additives, each additive is typically blended into the base oil in an amount which enables the additive to provide its desired function. Representative effective amounts of such additives, when used in crankcase lubricants, are listed below. All the values listed are stated as mass percent active ingredient. ______________________________________ Mass % Mass %Additive (Broad) (Preferred)______________________________________Ashless Dispersant 0.1-3 1-3Metal Detergents 0.1-15 0.2-9Corrosion Inhibitor 0-5 0-1.5Metal Dihydrocarbyl Dithiophosphate 0.1-6 0.1-4Anti-oxidant 0-5 0.01-2Pour Point Depressant 0.01-5 0.01-1.5Anti-Foaming Agent 0-5 0.001-0.15Supplemental Anti-wear Agents 0-0.5 0-0.2Friction Modifier 0-5 0-1.5Viscosity Modifier 0.01-10 0.25-3Low Saturate Base Oil Balance Balance______________________________________ In a preferred embodiment of the invention the oil comprises not more than 2 mass % of ashless dispersant and preferably does not contain monofunctional viscosity modifier. The components may be incorporated into a base oil in any convenient way. Thus, each of the components can be added directly to the oil by dispersing or dissolving it in the oil at the desired level of concentration. Such blending may occur at ambient temperature or at an elevated temperature. Preferably all the additives except for the viscosity modifier and the pour point depressant are blended into a concentrate or additive package described herein as the detergent inhibitor package, that is subsequently blended into basestock to make finished lubricant. Use of such concentrates is conventional. The concentrate will typically be formulated to contain the additive(s) in proper amounts to provide the desired concentration in the final formulation when the concentrate is combined with a predetermined amount of base lubricant. Preferably the detergent inhibitor package is made in accordance with the method described in U.S. Pat. No. 4,938,880. That patent describes making a premix of ashless dispersant and metal detergents that is pre-blended at a temperature of at least about 1000° C. Thereafter the pre-mix is cooled to at least 85° C. and the additional components are added. The final formulations may employ from 2 to 18 mass % and preferably 4 to 15 mass % of the concentrate or additive package (including any diluent or solvent contained in individual additives) with the remainder being viscosity modifier (in an appropriate amount to give the desired viscometrics) and base oil. The invention will now be described by of illustration only with reference to the following examples. EXAMPLE 1 An SAE 15W-40 oil of the invention prepared from a basestock of 64 mass % saturates was tested in the Sequence VE engine test, using a detergent inhibitor package with a reduced amount of ashless dispersant such that the level of active ingredient of the ashless dispersant is approximately 1.75 mass %. At a treat rate of 9.5 mass % of the preferred multifunctional viscosity modifier as described in WO-A-94/13763, without any monofunctional viscosity modifier a passing engine test result was obtained. Details of the oil and test result are set out in the Table below. ______________________________________Example 1______________________________________Basestock (mass %) 56.5% BP 150ME 24.0% BP 500ME Total saturates 64%Viscosity Modifier (mass %) 9.5% PARATONE 8500.sup.1Additive Package (mass %) 10.0% additive package.sup.2Sequence VE Engine Test ResultsSludge Rating (pass = 9.0 for API SH 9.1quality level)Varnish Rating (pass = 5.0 for API SH 6.0quality level)Cam Lobe Wear (pass = 5.0 for API SH 3.1quality level)______________________________________ Footnotes: .sup.1 multifunctional viscosity modifier according to WOA-94/13763 commercially available from Exxon Chemical Company and comprising an oil solution of a blend of derivatised polymers, with a polymer content of 10.2 mass %; .sup.2 a detergent inhibitor package comprising ashless dispersant, metalcontaining detergents, antioxidant, antiwear additive, antifoam additive, demulsifier, friction modifier and seal comparability aid.

Multigrade lubricating oils for use in lubricating internal combustion engines, using basestocks with low levels (<75 mass %) of saturated hydrocarbons, comprise less than 3 mass % of ashless dispersant derived from a polymer of number average molecular weight (Mn) of not greater than 5000, and a viscosity modifier package to give the desired viscometrics comprising at least one multifunctional viscosity modifier. These oils meet stringent engine performance requirements and specifically give adequate varnish inhibition without very high treat levels of dispersants and/or use of specific detergent systems so avoiding problems of oxidation stability, compatibility and engine performance debits.

This invention relates generally to biochemical and biological effects of nonionic nucleic acid methylphosphonates, and more particularly to nonionic nucleic acid alkyl and aryl methylphosphonates and processes for the manufacture and use thereof. Prior to the present invention, studies on nucleic acid analogs and derivatives possessing modified internucleoside linkages have made important contributions to understanding nucleic acid conformation in solution and have provided materials for various biochemical and biological studies. Recent studies have been made on the physical, biochemical and biological properties of one class of nonionic nucleic acid derivative, the oligonucleotide alkyl phosphotriesters. The physical properties of dinucleotide methyl and ethyl phosphotriesters have been studied by ultraviolet, circular dichroism, infrared and proton nuclear magnetic resonance spectroscopy. The interaction of deoxyribooligonucleotide ethyl phosphotriesters with sequences complementary to the amino acid accepting stem and anticodon region of transfer RNA have been characterized and their inhibitory effects on in vitro aminoacylation have been studied. More recently, the inhibitory effect of a 2'-O-methylribooligonucleotide triester, G p m (Et) G p m (Et)U, on cellular protein synthesis and growth of mammalian cells in culture has been reported. In addition, selective binding of an octathymidylate ethyl phosphotriester, [Tp(Et)] 7 T to polydeoxyadenylic acid has been extensively investigated. An object of this invention is to provide nonionic nucleic acid alkyl or aryl phosphonates analogs. Still another object of this invention is to teach the preparation of nonionic nucleic acid alkyl or aryl phosphonate analogs by several novel synthetic processes or methods. To provide nonionic nucleic acid alkyl or aryl phosphonate analogs for interacting with complementary cellular or viral nucleic acids with the objective of controlling or regulating the function or expression of the cellular or viral nucleic acids, is still another object of this invention. To provide a heptadeoxyribonucleoside methyl phosphonate with a base sequence which is complementary to the 3'-terminus of bacterial 16S ribosonal ribonucleic-acid with objective of preventing bacterial protein synthesis, is a further object of this invention. Even another object of this invention is to provide alkyl and aryl phosphonate nucleic acid analogs comprising at least two nucleosides, a base, and an alkyl or aryl phosphonate group, with the nucleosides being linked together to form alkyl or aryl phosphonate nucleic acid analogs for the purpose of controlling or regulating the function or expression of cellular or viral nucleic acids. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 shows the general molecular structure of an oligonucleoside alkyl or aryl phosphonate; FIG. 2 illustrates the general molecular structure of an oligonucleoside alkyl or aryl phosphonate; FIG. 3 illustrates the molecular structure of a heptadeoxyribonucleoside methylphosphonate; FIG. 4 is a schematic of a process for the synthesis of a dinucleoside alkyl or arylphosphonate; FIG. 5 is a schematic of a process for the synthesis of a dinucleoside alkyl or arylphosphonate; FIG. 6 is a schematic of a process for the synthesis of a dinucleoside alkyl or arylphosphonate; FIG. 7 is a schematic of a process for the synthesis of a dinucleoside alkyl or arylphosphonate; FIG. 8 illustrates the process for the synthesis of a dinucleoside alkyl or arylphosphonate; FIG. 9 illustrates the process for the synthesis of a dinucleoside alkyl or arylphosphonate; FIG. 10 shows a process for the synthesis of a dinucleoside alkyl or arylphosphonate; FIG. 11 illustrates the process for synthesis of a dinucleoside alkyl or arylphosphonate; FIG. 12 is spectral diagrams of the 360 MHz 1 H nmr spectra of (a) d-ApA) 1 ; (b) d-ApA) 2 ; (c) d-TpT) 1 ; and (d) d-TpT) 2 at 25° C. in D 2 O containing 1 mM ethylenediaminetetracetate-10 mM sodium phosphate, pH 7.0, with the tentative chemical shift assignments appear above each dimer; FIG. 13 is spectral diagrams of circular dichroism spectra of (a) dApA) 1 (--); dApA) 2 ( - - - ); (b) dApT) 1 (--), dApT 2 ( - - - ); (c) dTpA) 1 ( - - - ), dTpA) 2 (--); and (d) dTpT in 10 mM Tris.HCl, 10 mM MgCl 2 , pH 7.5 at 27° C.; FIG. 14 is a sketch showing mixing experiment between polyuridylic acid and dApA) 1 (O) or dApA) 2 (o) in 10 mM Tris, 10 mM MgCl 2 , pH 7.5 at 0° C., with the total nucleotide concentration is 1×10 -4 M. FIG. 15 is a sketch showing melting curves of poly U+dApA) 1 (O) and poly U+dApA) 2 (o) in 10 mM Tris, 10 mM MgCl 2 , pH 7.5, with the stoichiometry of each complex is 2U:1A and the total nucleotide concentration is 5×10 -5 M; FIG. 16 is molecular diagrams of the diastereoisomers of dideoxyribonucleoside methyl phosphonates; FIG. 17 is a schematic of a synthetic of this invention for preparation of the oligonucleoside methyl phosphonate; and FIG. 18 illustrates the transport of (O) 100 μM d-GpGp-[ 3 H]-T and (o) 100 μM d-Tp) 8 -[ 3 H]-T into transformed Syrian hamster fibroblasts growing in monolayer at 37° C. DESCRIPTION OF THE PREFERRED EMBODIMENTS The general structure of the oligonucleoside alkyl or aryl phosphonates are shown in FIGS. 1 and 2. FIG. 1 shows the general structure of an oligonucleoside alkyl or aryl phosphonate. The nucleoside units which consists of a base (B) comprising adenine, thymine, guanine, cytosine, uracil or hypoxanthine and a sugar where R' can be hydrogen, hydroxyl, O-alkyl or O-aryl or O-halogeno are linked in a 3'-5' manner by a phosphosphonate group where R can be alkyl or aryl. The configuration of the alkyl or aryl phosphonate group is S. FIG. 2 shows the general structure of an oligonucleoside alkyl or aryl phosphonate. The nucleoside units which consists of a base (B) comprising adenine, thymine, guanine, cytosine, uracil or hypoxanthine and a sugar where R' can be hydrogen, hydroxyl, O-alkyl or O-aryl or O-halogeno are linked in a 3'→5' manner by a phosphosphonate group where R can be alkyl or aryl. The configuration of the alkyl or aryl phosphonate group is R. A specific example of an oligodeoxyribonucleoside methylphosphonate is shown in FIG. 3. In FIG. 3, there is shown the structure of a heptadeoxyribonucleoside methylphosphonate. The deoxyribonucleoside units which occur in the order deoxadenosine-deoxyguanosine-deoxyquanosine-deoxyadenosine-deoxyguanosine-deoxyquanosinethymidine are linked in a 3'→5' manner by methylphosphonate groups. The configuraitons of the methylphosphonate groups are not specified. The following description will now be given of the synthesis of a particular series of oligodeoxyribonucleoside methyl phosphonates and their physical properties. The synthetic procedure resulted in the separation of two diastereoisomers of each dimer analog. This synthetic scheme also allowed the preparation of analogs containing a 13 C-enriched phosphonate methyl group. The influence of backbone configuration on overall dimer conformation was studied by ultraviolet, circular dichroism and 1 H, 13 C and 31 P nuclear magnetic resonance techniques and the results were compared to the conformations of their parent dideoxyribonucleoside monophosphates. Furthermore, the effects of backbone configuration and the removal of the negative charge on the interaction of deoxyadenosine-containing dimers with polyuridylic acid and polythymidylic acid were assessed. Several processes or methods will now be described for synthesis of the various materials of this invention. FIG. 4 shows the synthesis of a dinucleoside alkyl or arylphosphonate. This process consists of esterification of a 5'-O-protected nucleoside having 3'-hydroxyl group with an alkyl or arylphosphonic acid in the presence of an activating agent to form a 5'-O-protected nucleoside-3'-O-alkyl or arylphosphonate. The latter compound is then esterified with a 3'-O-protected nucleoside having a 5'-hydroxyl group in the presence of an activating agent to form a fully protected dinucleoside alkyl or arylphosphonate. The protecting groups are then removed from the fully protected dinucleoside alkyl or arylphosphonate to form the dinucleoside alkyl or arylphosphonate. Referring now to FIG. 5, there is illustated the synthesis of a dinucleoside alkyl or arylphosphonate. This process consists of esterification of a 3'-O-protected nucleoside having a 5'-hydroxyl group with an alkyl or arylphosphonic acid in the presence of an activating agent to form a 3'-O-protected nucleoside-5'-O-alkyl or arylphosphonate. The latter compound is then esterified with a 5'-O-protected nucleoside having a 3'-hydroxyl group in the presence of an activating agent to form a fully protected dinucleoside alkyl or arylphosphonate. The protecting groups are then removed from the fully protected dinucleoside alkyl or arylphosphonate to form the dinucleoside alkyl or arylphosphonate. FIG. 6 shows the synthesis of a dinucleoside alkyl or arylphosphonate. This process consists of esterification of a 5'-O-protected nucleoside having a 3'-hydroxyl group with disubstituted alkyl or arylphosphonate. The substitutents (R"") can be chloride, imidazolide, triazolide or tetrazolide. A 5'-O-protected nucleoside-3'-O-monosubstituted alkyl or arylphosphonate is formed in this reaction. This compound is then esterified with a 3'-protected nucleoside having a 5'-hydroxyl group to give the fully protected dinucleoside alkyl or arylphosphonate. The protecting groups are then removed from the fully protected dinucleoside alkyl or arylphosphonate to form the dinucleoside alkyl or arylphosphonate. The synthesis of a dinucleoside alkyl or arylphosphonate is shown in FIG. 7. This process consists in the esterification of a 3'-O-protected nucleoside having a 5'-hydroxyl group with disubstituted alkyl or arylphosphonate. The substituents (R"") can be chloride, imidazolide, triazolide or tetrazolide. A 3'-O-protected nucleoside-5'-O-monosubstituted alkyl or arylphosphonate is formed in this reaction. This compound is then esterified with a 5'-protected nucleoside having a 3'-hydroxyl group to give the fully protected dinucleoside alkyl or arylphosphonate. The protecting groups are then removed from the fully protected dinucleoside alkyl or arylphosphonate to form the dinucleoside alkyl or arylphosphonate. FIG. 8 illustrates the synthesis of a dinucleoside alkyl or arylphosphonate. In this process a 5'-O-protected nucleoside having a 3'-hydroxyl group is reacted with a dichloro alkyl or arylphosphine to form a 5'-O-protected nucleoside-3'-O-monochloro alkyl or arylphosphonate. This compound is then reacted with a 3'-O-protected nucleoside having a 5'-hydroxyl group to give a fully protected dinucleoside alkyl or arylphosphonate. The latter compound is then oxidized to form a fully protected dinucleoside alkyl or arylphosphonate. The protecting groups are then removed from the fully protected dinucleoside alkyl or arylphosphonate to give the dinucleoside alkyl or arylphosphonate. FIG. 9 shows the synthesis of a dinucleoside alkyl or arylphosphonate. In this process, a 3'-O-protected nucleoside having a 5'-hydroxyl group is first reacted with a dichloro alkyl or arylphosphine to form a 3'-O-protected nucleoside-5'-O-monochloro alkyl or arylphosphonate. This compound is next reacted with a 5'-O-protected nucleoside having a 3'-hydroxyl group to give a fully protected dinucleoside alkyl or arylphosphinate. The latter compound is then oxidized to form a fully protected dinucleoside alkyl or arylphosphonate. The protecting groups are then removed from the fully protected dinucleoside alkyl or arylphosphonate to give the dinucleoside alkyl or arylphosphonate. FIG. 10 shows the synthesis of a dinucleoside alkyl or arylphosphonate. In this process, a 5'-O-protected nucleoside having a 3'-hydroxyl group is first reacted with ethyldichlorophosphite to give a 5'-O-protected nucleoside-3'-O-ethyl monochlorophosphite. This compound is then reacted with a 3'-O-protected nucleoside having a 5'-hydroxyl group to give a fully protected dinucleoside ethylphosphite. The latter compound is then reacted with an alkyl or aryliodide to give a fully protected dinucleoside alkyl or arylphosphonate. The protecting groups are then removed from the fully protected dinucleoside alkyl or arylphosphonate to give the dinucleoside alkyl or arylphosphonate. FIG. 11 shows the synthesis of a dinucleoside alkyl or arylphosphonate. In this process, a 3'-O-protected nucleoside having a 5'-hydroxyl group is reacted with ethyldichlorophosphite to give a 3'-O-protected nucleoside-5'-O-ethyl monochlorophosphite. This compound is reacted with a 5'-O-protected nucleoside having a 3'-hydroxyl group to give a fully protected dinucleoside ethylphosphite. The latter compound is then reacted with an alkyl or aryliodide to give a fully protected dinucleoside alkyl or arylphosphonate. The protecting groups are then removed from the fully protected dinucleoside alkyl or arylphosphonate to give the dinucleoside alkyl or arylphosphonate. In previous processes (FIGS. 4 to 11), there was described several methods for joining nucleosides to form dinucleoside alkyl or aryl phosphonates. The chain length can be extended from at least 2 to a finite number of greater than 2, for example up to 20. For rapid and efficient synthesis oligonucleoside alkyl or aryl phosphonates can be joined by processes analogous to there previously described for the preparation of dinucleoside alkyl or aryl phosphonates. Materials and Methods Thymidine and 2'-deoxyadenosine were checked for purity by paper chromatography before use. 5'-Mono-p-methoxytrityl thymidine, 5'-di-p-methoxy-trityl-N-benzoyldeoxyadenosine, 3'-O-acetylthymidine and 3'-O-acetyl-N-benzoyldeoxyadenosine were prepared according to published procedures. Diethyl [ 13 C]-methyl phosphonate was prepared by reaction of [ 13 C]-methyl iodide (96% enriched) with triethyl phosphate followed by vacuum distillation of the product (b.p. 64°-66°/2 mm/Hg. The pyridinium salt of methyl phosphonic acid was prepared by hydrolysis of dimethyl methylphosphonate or diethyl [ 13 C]-methyl phosphonate in 4N hydrochloric acid followed by isolation of the product as the barium salt. The barium salt was converted to the pyridinium salt by passage through a Dowex 50X pyridinium ion-exchange column. Mesitylenesulfonyl chloride was treated with activated charcoal and recrystallized from pentane immediately before use. 1-H-Tetrazole was then prepared. All solvents and reagents were purified. Silica gel column chromatography was performed using Baker 3405 silica gel (60-200 mesh). Thin layer silica gel chromatography and thin layer cellulose chromatography were done. Paper chromatography was carried out on Whatmann 3 MM paper using the following solvent systems: solvent A, 2-propanol-concd. ammonium hydroxide-water (7:1:2 v/v); solvent C, 1M ammonium acetate-95% ethanol (3:7, v/v); solvent F, 1-propanol-concd. ammonium hydroxide-water (50:10:35, v/v) or solvent I, 2-propanol-water (7:3, v/v). High pressure liquid chromatography was performed using columns (2.1 mm×1 m) packed with DuPont Permaphase ODS reverse phase material. Linear gradients (40 ml) from 0% to 75% methanol in water were used at a flow rate of 1 ml/min. The HPLC mobility refers to the percentage of methanol in water required to elute the compound from the column. For reactions carried out in pyridine, the reactants were dried by repeated evaporation with anhydrous pyridine and were then dissolved in anhydrous pyridine. Unless otherwise noted, all reactions and operations were performed at room temperature. Preparation of Mesitylenesulfonyl Tetrazolide Although the preparation of MST has been published, a modified procedure was used. A solution of 1-H-tetrazole (3.36 g: 48 mmoles) in 40 ml of dry methylene chloride containing 5.6 ml (40 mmoles) of triethylamine was added dropwise to 40 ml of anhydrous ether containing mesitylenesulfonyl chloride (8.76 g, 40 mmoles) at room temperature. After 2 hrs, the precipitated triethylammonium chloride was removed by filtration, and washed with 50 ml of methylene chloride-ethyl ether (1:1 v/v). The filtrate was cooled to 0° C. and pentane was added until the solution became cloudy. This procedure was repeated until a total volume of 25 ml had been added over a 4 hour period. After storage overnight at 4°, the resulting white solid was removed by filtration on a sintered glass filter. The solid was dissolved in 500 ml of anhydrous ethyl ether. The solution was then filtered to remove a small amount of insoluble material. The filtrate was evaporated to dryness and the resulting mestilylene sulfonyl tetrazolide (4.46 g) was obtained in 44% yield. The MST was pure as indicated by silica gel tlc, Rf (C 6 H 6 ) 0.11 (m.p. 109°-110° C.) and was stored in a desicator at -80° C. Under these storage conditions, the MST maintained its condensing activity for at least one month. Preparation of [MeOTr]TpCE [MeOTr]T (10.3 g, 20 mmoles), the pyridinium salt of methyl phosphonic acid (40 mmoles) and Dowex pyridinium resin (0.5 g) were treated with dicyclohexyl-carbodiimide (41.2 g, 200 mmoles) in 100 ml of pyridine at 37° C. for 3 days. The resulting [MeOTr]Tp, Rf silica gel tlc 0.00 (EtOAc/THF 1:1), appeared to be formed in approximately 86% yield as determined by HPLC. The material is eluted from the HPLC column with 37% methanol/water. Hydracrylonitrile (100 ml) was added to the reaction mixture, which was kept at 37° C. for 2 days. Aqueous pyridine (200 ml) was then added and the resulting dicyclohexylurea was removed by filtration. The filtrate was evaporated, dissolved in 250 ml of ethyl acetate and the solution extracted with 3 (250 ml) portions of water. The ethyl acetate solution was dried over anhydrous sodium sulfate. After filtration and evaporation, the mixture was chromatographed on a silica gel column (5.4×37 cm) which was eluted with ether (IL), ethyl acetate (1.2 L) and tetrahydrofuran (1.6 L). Pure [MeOTr]TpCE (7.5 g) was isolated in 55% yield after precipitation from tetrahydrofuran by addition of hexane. The monomer has Rf values of 0.32 (EtOAc/THF 1:1) and 0.66 (20% MeOH-CHCl 3 ) on silica gel tlc and is eluted from the HPLC column with 54% methanol/water. The UV spectrum gave λ max 267 nm, sh 230 nm, λ min 250 nm, ε260/ε280 1.44 in 95% ethanol. Preparation of [MeOTr]Tp Pyridinium Salt [MeOTr]TpCE (2.17 g, 3.36 mmoles) was treated with 16.8 ml of 1N sodium hydroxide in a solution containing 126 ml of dioxane and 25 ml of water for 15 min. The solution was neutralized by addition of Dowex 50X pyridinium resin. After filtration, the solution was evaporated and the residue was rendered anhydrous by evaporation with pyridine. The [MeOTr]Tp (1.90 g; 2.83 mmoles) was obtained in 84% yield after precipitation from pyridine by addition to anhydrous ether. The material has Rf values of 0.00 (EtOAc/THF 1:1) and 0.04 (20% MeOH/CHCl 3 ) on silica gel tlc. The UV spectrum gave λ max 267 nm, sh 230 nm, λ min 250 nm, ε230/ε267 1.68, ε260/ε280 1.44 in 95% ethanol. The monomethoxytrityl group was removed from 70 mg (0.1 mmole) of d-[(MeO)Tr]Tp by treatment with 80% aqueous acetic acid. The resulting Tp (874 A 267 units, 0.095 mmole) was isolated in 95% yield by chromatography on a DEAE Sephadex A25 column (3×8.5 cm) using a linear gradient of ammonium bicarbonate (0.01M to 0.20M, 500 ml). The monomer has the following Rf values on cellulose TLC:0.41 (solvent A), 0.77 (solvent C) and 0.69 (solvent F). The UV spectrum gave λ max 267 nm, λ min 235 nm in water pH 7.0. The pmr spectrum was consistent with the structure of the monomer. Preparation of TpCE [MeOTr]TpCE (3.26 g, 5.04 mmoles) dissolved in 20 ml of methanol was treated with 80 ml of 80% acetic acid solution for 5 hr. at 37° C. The solvents were removed by evaporation and the residue was evaporated repeatedly with 50% toluenetetrahydrofuran to remove the acetic acid. TpCE (1.80 g, 4.8 mmoles) was obtained in 96% yield after precipitation from tetrahydrofuran (10 ml) by addition of hexane (200 ml). The material has Rf values of 0.08 (EtOAc/THF, 1:1) and 0.16 (15% MeOH/CHCl 3 ) on silica gel tlc. The UV spectrum gave λ max 265 nm, λ min 233 nm, ε260/ε280 1.61 in absolute ethanol. Preparation of d-[(MeO) 2 Tr]bzApCE A solution containing d-[MeO) 2 Tr]bzA (10.5 g; 16 mmoles), methyl phosphonic acid (32 mmoles) and Dowex 50X pyridinium resin (0.5 g) in 80 ml of anhydrous pyridine was treated with dicyclohexylcarbodiimide (25 g; 121 mmoles) for 3 days at 37° C. Examination of the reaction mixture by HPLC showed essentially quantitative conversion of d-[(MeO) 2 Tr]bzA to d-[(MeO) 2 Tr]bzAp, which has HPLC retention time of 22.8 min. The reaction mixture was treated with 80 ml of hydracrylonitrile for 2 days at 37° C. After filtration and evaporation of the solvents, the residue was dissolved in 200 ml of ethyl acetate and the solution was extracted with three (200 ml) portions of water. The ethyl acetate solution was dried over anhydrous sodium sulfate, concentrated to 50 ml and chromatographed on a silica gel column (5.4×37 cm). The column was eluted with ether (1.5 L), ethyl acetate (1.5 L) and tetrahydrofuran (1.5 L). The resulting d-[(MeO) 2 Tr]bzApCE weighed 5.4 g (6.84 mmoles, 43%) after precipitation from tetrahydrofuran (100 ml) with hexane (500 ml). The material elutes from the HPLC column with 68% methanol/water, and has silica gel TLC Rf values of 0.13 (EtOAc/THF, 1:1 v/v) and 0.27 (THF). The UV spectrum shows λ max 279 nm and 234 nm, λ min 258 nm and 223 nm; ε234/ε279 1.44, ε260/ε280 0.67, in 95% ethanol. Preparation of d-[(MeO) 2 Tr]bzAp A solution containing d-[(MeO) 2 Tr]bzApCE (3.79 g, 4.8 mmoles) in 180 ml of dioxane and 36 ml of water was treated with 24 ml of 1N sodium hydroxide for 7 min. The solution was neutralized with Dowex 50X pyridinium resin and then was passed through a Dowex 50X pyridinium ion exchange column (3×30 cm). The eluate was evaporated and the residue was rendered anhydrous by evaporation with pyridine. The resulting d-[(MeO) 2 Tr]bzAp (2.9 g; 3.56 mmoles) was obtained in 74% yield after precipitation from anhydrous ether. The material has Rf values of 0.00 (THF) and 0.36 (50% MeOH/CHCl 3 ) on silica gel tlc. The UV spectrum showed λ max 280 nm and 233 nm, λ min 255 nm and 225 nm ε233/ε280 1.37, ε260/ε280 0.67 in 95% ethanol. The protecting groups were removed from a small sample of d-[(MeO) 2 Tr]bzAp (80 mg, 0.1 mmole) by sequential treatment with concentrated ammonium hydroxide in pyridine and 80% acetic acid. The monomer dAp (1400 A 260 units, 0.09 mmole) was isolated by chromatography on a DEAE Sephadex A-25 column (3×8.5 cm) using a linear gradient of ammonium bicarbonate (0.01 to 0.2M, 600 ml). The monomer has the following Rf values on cellulose TLC:0.45 (solvent A), 0.56 (solvent C) 0.67 (solvent F) and 0.44 (solvent I). The UV spectrum showed λ max 259 nm, λ min 227 nm 260/280 6.13, in water pH 7.0. The pmr spectrum was consistent with the structure of the monomer. Preparation of d-bzApCE A solution of d-[(MeO) 2 Tr]bzApCE (1.58 g, 2 mmoles) in 6.3 ml of methanol was treated with 25 ml of 80% acetic acid for 1.5 hrs. The solvents were evaporated and the residue was repeatedly evaporated with toluene and tetrahydrofuran to remove acetic acid. The residue was precipitated from 20 ml of tetrahydrofuran by dropwise addition to 250 ml of hexane to give 0.95 g (1.95 mmoles) of d-bzApCE in 98% yield. The material has Rf values of 0.09 (THF) and 0.25 (20% MeOH/CHCl 3 ) on silica gel TLC and is eluted from the HPLC column with 12% methanol/water. The UV spectrum shows λ max 280 nm, sh 233 nm, λ min 247 nm ε233/ε280 0.65 ε260/ε280 0.60, in 95% ethanol. Preparation of Dinucleoside Methyl Phosphonates The general procedure for the preparation of protected dinucleoside methyl phosphonates is given in this section. Table 1 shows the specific reaction conditions and yields for each dimer. The protected nucleoside 3'-methyl phosphonate and protected nucleoside or nucleoside 3'-methyl phosphonate cyanoethyl ester were dried by evaporation with anhydrous pyridine. The condensing agent was added and the reactants were taken up in anhydrous pyridine to give a 0.2M solution. After completion of the reaction as indicated by TLC and/or HPLC, an equal volume of water was added and the solution was kept at room temperature for 30 min. The solvents were then evaporated and the residue dissolved in ethyl acetate or chloroform. The organic solution was extracted with water and then dried over anhydrous sodium sulfate. After concentration, the organic solution was applied to a silica get column (3×28 cm for a 1 mmole scale reaction). The column was eluted with ethyl acetate, ethyl acetate/tetrahydrofuran (1:1 v/v) and tetrahydrofuran. The progress of the elution was monitored by silica get TLC. Dimers terminating with a 3'-acetyl group separated into their individual diastereoisomers on the column and were eluted as pure isomer 1, a mixture of isomer 1 and 2 and pure isomer 2. The dimers were isolated as white solids, by precipitation from tetrahydrofuran solution upon addition of hexane. The Rf values on silica get tlc, the mobilities on the HPLC column and the ultraviolet spectral characteristics of the protected dimers are given in Table 2. TABLE 1__________________________________________________________________________Preparation of Protected Dideoxyribonucleoside Methyl Phosphonates Condensing Agent.sup.(a) Reaction DimerMonomers (mmole) (mmole) Time (mmole) Yield__________________________________________________________________________d-[(MeO)Tr]Tp (0.20) DCC 3 days d-[(MeO)Tr]TpTOAc 16%+ d-TOAc (0.22) (0.73) 37° C. (0.031)d-[(MeO)Tr]Tp (2.40) MST 3 hrs. d-[(MeO)Tr]TpTpCE 55%+ d-TpCE (3.60) (9.60) r.t. (1.32)d-[(MeO).sub.2 Tr]bzAp (1.26) TPSCl 4 days d-[(MeO).sub.2 Tr]bzApbzAOAc 39%+ d-bzAOAc (1.50) (2.0) 37° C. (0.50)d-[(MeO).sub.2 Tr]bzA- (0.70) MST 4 hrs. d-[(MeO).sub.2 Tr]bzA-[.sup.13 C]--pbzAOAc 41%[.sup.13 C]--p + d-bzAOAc (1.05) (2.8) r.t. (0.29)d-[(MeO).sub.2 Tr]bzAp (0.85) MST 6 hrs. d-[(MeO).sub.2 Tr]bzApbzApCE 46%+ d-bzApCE (1.28) (4.0) r.t. (0.39)d-[(MeO)Tr]Tp (1.0) TPSCl 16 hrs. d-[(MeO)tr]TpbzAOAc 35%+ d-bzAOAc (1.0) (3.0) 37° C. (0.35)d-[(MeO).sub.2 Tr]bzAp (1.0) TPSCl 46 hrs. d-[(MeO).sub.2 Tr]bzApTOAc 38%+ d-TOAc (1.3) (1.5) 37° C. (0.38)__________________________________________________________________________ .sup.(a) DCC -- dicyclohexylcarbodiimide TPSCl -- triisopropylbenzenesulfonyl chloride MST -- mesitylenesulfonyl tetrazolide The base labile protecting groups were removed from the dimers by treatment with 50% concentrated ammonium hydroxidepyridine solution for 3 days at 4° C. Alternatively, the N-benzoyl protecting groups of dimers containing deoxyadenosine could be removed by treatment with 85% hydrazine hydrate in 20% acetic acid-pyridine buffer overnight at room temperature (Letsinger et al., 1968). This treatment also partially removed the 3'-O-acetyl group. The acetyl group was completely removed by further treatment with 50% concentrated ammonium hydroxide-pyridine solution for 2 hours at 4° C. After complete removal of solvents, the trityl protecting groups were removed by treatment with 80% acetic acid-methanol (8:2 v/v) solution at room temperature. The solvents were then removed and the dimers were chromatographed on Whatmann 3 MM paper using solvent A. The dimers were eluted from the paper with 50% aqueous ethanol. For dimers terminating with 3'-OH groups, the ethanol solutions were passed through small (0.5×1 cm) DEAD cellulose columns to remove trace impurities eluted from the paper chromatogram. Dimers terminating with 3'-methyl phosphonate groups were absorbed to small DEAE cellulose columns and then eluted with 0.5M ammonium bicarbonate solution. The dimers were stored as standard solutions in 50% ethanol at 0° C., and were found to be complete stable under these conditions for at least 9 months. For physical and nmr studies, aliquots containing the required amount of dimer were evaporated to remove the ethanol and then lyophilized from water or D 2 O before use. The Rf values and UV spectral characteristics of the dimers are given in Table 3. The pmr spectra and tentative chemical shift assignments of the two diastereoisomers of d-ApA and d-TpT are shown in FIG. 12. TABLE 2__________________________________________________________________________Chromatographic Mobilities and Ultraviolet Spectral Propertiesof Protected Dideoxyribonucleoside Methyl Phosphonates ##STR1## HPLC Mobility.sup.(c) 10% MeOH/ 15% MeOH/ (% methanol Ultraviolet Spectral Properties.sup.(b)Dimer THF CHCl.sub.3 CHCl.sub.3 water) λ max (nm) λ min__________________________________________________________________________ (nm)d-[(MeO)Tr]TpTOAc 0.44 0.63 0.53 -- -- 267sh235 246 ##STR2##d-[(MeO)Tr]TpTpCE -- -- 0.28 0.22 -- 265sh235 245 ##STR3##d-[(MeO).sub.2 Tr]bzApbzAOAc 0.34 0.29 -- 0.51 0.43 68% 281sh233 256 ##STR4##d-[(MeO).sub.2 Tr]bzApbzApCE 0.09 0.32 0.28 0.39 0.35 -- 281sh230 256 ##STR5##d-[(MeO)Tr]TpbzAOAc 0.47 0.24 0.19 0.59 62% 276;230, sh260 247;227 ##STR6##d-[(MeO).sub.2 Tr]bzApTOAc 0.49 0.41 -- 0.52 0.48 66% 277;235, sh263 255,227 ##STR7##__________________________________________________________________________ .sup.(a) Two Rf values refer to the mobilities of the individual diastereoisomers. .sup.(b) Ultraviolet spectra were measured in 95% ethanol at room temperature. .sup.(c) Percentage of methanol in water required to elute compound from HPLC column (DuPont Phermaphase ODS) TABLE 3______________________________________Chromatographic Mobilities and Ultraviolet SpectralProperties of Dideoxyribonucleoside Methyl Phosphonates Ultraviolet SpectralMobility Paper Properties.sup.(a) Dimer Chromatography Rf(A)Rf(C)Rf(I) λ max (nm) λ min (nm) ##STR8##______________________________________d-TpT -- 0.67 0.73 267 234 1.91d-ApA -- 0.48 0.62 258 223 6.65d-ApAp 0.30 -- 0.38 258 227 4.52d-TpA 0.42 0.57 0.63 262 230 2.80d-ApT 0.33 0.55 0.61 261 233 3.22______________________________________ TABLE 4______________________________________Hypochromicity of Dideoxyadenosine Methyl PhosphonateAnalogsCompound ε(molar).sup.(a) % Hypochromicity______________________________________dpA 15.3 × 10.sup.3 --d-ApA 12.7 × 10.sup.3 17%(d-ApA).sub.1 13.7 × 10.sup.3 11.0%(d-ApA).sub.2 14.3 × 10.sup.3 7.1%(d-ApAp).sub.1 13.0 × 10.sup.3 13.3%(d-ApAp).sub.2 13.3 × 10.sup.3 11.3%______________________________________ .sup.(a) Measured in 1 mM Tris HCl pH 7.4 at 27° C. Similar pmr spectra were obtained for d-ApT and d-TpA (data not shown). The spectra are consistent with the structures of the dimers. The complete characterization of all these dimers by pmr spectroscopy will be described in a subsequent paper (Kan et al., manuscript in preparation). Physical Studies and Interaction with Polynucleotides Ultraviolet and circular dichroism spectra were recorded respectively on a Cary 15 spectrophotometer and a Cary 60 spectropolarimeter with CD attachment. The continuous variation experiments, melting experiments and circular dichroism experiments were carried out as previously described. The molar extinction coefficient of poly U is 9.2×10 3 (265 nm) and poly dT is 8.52×10 3 (264 nm). The molar extinction coefficients of the dideoxyadenosine methyl phosphonates were determined by comparing the absorption of a solution of the dimer at pH 7.4 with the absorption of the same solution at pH 1.0. The dimer extinction coefficient was then calculated from the observed hyperchromicity of the dimer at pH 1.0 using an extinction coefficient for deoxyadenosine at pH 1.0 of 14.1×10 3 . Preparation of Dinucleoside Methyl Phosphonates The synthetic route used to prepare the dinucleoside methyl phosphonates has previously been described. 5'-Mono-p-methoxytrityl thymidine and 5'-di-p-methoxytrityl-N-benzoyl deoxyadenosine were converted to the corresponding 3'-methyl phophonate β-cyanoethyl esters (2) by sequential reaction of (1) with methyl phosphonic acid and β-cyanoethanol in the presence of dicyclohexylcarbodiimide. The preparation of nucleoside 5'-methyl phosphonates by reaction of a suitably protected nucleoside with methyl phosphonic acid has been previously known. Direct conversion of the protected nucleoside-3'-methyl phosphonate to its β-cyanoethyl ester allows purification of this intermediate on a large scale by silica gel column chromatography, thus avoiding the use of ion exchange chromatography. They trityl or β-cyanoethyl protecting groups can be selectively removed from 2 by treatment with either 80% acetic acid or 0.1N sodium hydroxide solution, respectively, at room temperature. Protected nucleoside-3'-methyl phosphonate (4) was condensed with either 3'-O-acetyl thymidine or 3'-O-acetyl-N-benzoyldeoxyadenosine to give fully protected dinucleoside methyl phosphonate 6. Alternatively, 4 was condensed with the β-cyanoethyl ester of thymidine-3'-methyl phosphonate or N-benzoyldeoxyadenosine-3'-methyl phosphonate to give 7. The condensing agents used in these reactions were dicyclohexylcarbodiimide, triisopropylbenzenesulfonyl chloride or mesitylenesulfonyl tetrazolide. The reaction conditions and yields are given in Table 1. The fully protected dimers were readily purified by silica gel column chromatography. For dimers terminating with 3'-O-acetyl groups, the two diastereoisomers were sufficiently separated on the silica gel column that fractions containing each pure diastereoisomer were obtained. These isomers were designated isomer 1 and isomer 2 in reference to their order of elution from the column. The diastereoisomers were generally formed in a 4:6 ratio of isomer 1 to isomer 2. Alternatively, the diastereoisomers could be obtained in pure form by thick layer chromatography on silica gel plates. The dimers terminating in a 3'-O-β-cyanoethyl methyl phosphosphonate group (7) consists of four diastereoisomers, although only two separate bands were observed on silica gel thin layer chromatography (see Table 2). For the deoxyadenosine-containing dimer, these two bands turned out to the two isomers with opposite configuration (axial and equatorial, see FIG. 16) about the methyl phosphonyl internucleoside linkage. Removal of the protecting groups from 6 and 7 was accomplished by sequential treatment with concentrated ammonium hydroxide in pyridine for 3 days at 4° C. followed by treatment with 80% acetic acid. In the case of the dideoxyadenosine methyl phosphonates, some hydrolysis of the phosphonate linkage was noted when the ammonium hydroxide treatment was carried out at room temperature. However, the hydrolysis was suppressed at low temperature. Alternatively, the N-benzoyl protecting groups of these dimers could be removed by treatment with hydrazine hydrate. The dimers were then purified by paper chromatography. The individual diastereoisomers of each deprotected dimer had the same chromatographic mobilities on paper chromatography in all solvent systems tested (See Table 3). Ultraviolet and Hypochromicity Measurements The ultraviolet spectral properties of the dinucleoside methyl phosphonates are recorded in Table 3. Qualitatively, the spectra are similar to those of (3'-5')-linked dinucleoside monophosphates. The spectra of the individual diastereoisomers are qualitatively similar to each other. Hypochromicity measurements for the dideoxyadenosine methyl phosphonates were carried out in water at pH 7.4 and are shown in Table 4. The percent hypochromicity of the methyl phosphonate dimers is from 4% to 10% lower than the percent hypochromicity of d-ApA. Each diastereoisomer has an unique molar extinction coefficient. The hypochromicity of isomer 1, the isomer eluted first from the silica gel column, is greater that that of isomer 2, reflecting differences in the extent of base-base overlap in these dimers. Circular Dichroism Spectra Differences in the extent and mode of base stacking interactions are observed for individual diastereoisomers within a given dimer sequence as reflected by the CD spectra of the dimers. The profile of the CD spectrum of d-ApA) 1 (FIG. 13a) is qualitatively similar to that of the parent dinucleoside monophosphate, d-ApA (Miller et al., 1971). However, the magnitudes of the molecular ellipticity (θ) at 267 nm and 270 nm of d-ApA) 1 are approximately one-half of those found for d-ApA. A very dramatic difference in the CD spectrum of d-ApA) 2 is observed. Only negative [θ] is found at 250 nm and the amplitude of the molecular ellipticity is approximately three-fold less than that of d-ApA) 1 . Similar results were observed for d-ApAp) 1 and d-ApAp) 2 (data not shown). In the case of d-ApT (FIG. 13b), the profiles of the CD spectra of both isomers 1 and 2 are qualitatively similar to that of d-ApT (Cantor et al., 1970). However, the magnitude of the ellipticity of the peak (272 nm) and trough (253 nm) of the dinucleoside methyl phosphonate are less than those in the dinucleoside monophosphate. For d-ApT 1 , the peak is reduced 1.8-fold and the trough is reduced 1.3-fold compared to d-ApT while for d-ApT) 2 the reductions are 8.1 and 5.0-fold. The CD spectra of d-TpA) 1 and d-TpA) 2 (FIG. 13c) show differences in both the magnitude of the molecular ellipticity and in the position of the positive and negative bands. Isomer 2 has a CD spectrum which is virtually identical to that observed for d-TpAp (Cantor et al., 1970). Isomer 1, on the other hand, has a lower magnitude of the (θ) value, while the positions of the peak and trough are shifted to shorter wavelengths. The CD results for dApT) 1 and 2 and dTpA) 1 and 2 are qualitatively similar to those obtained by others on the CD spectra of the separated disastereoisomers of the dinucleoside ethyl phosphotriesters, dAp(Et)T and dTp(Et)A. In the case of these triesters, one isomer has a spectrum which is almost identical to that of the corresponding dinucleoside monophosphate. The other isomer shows significant reductions in the magnitudes of both the positive and negative CD absorbtion bands. It is not possible at this time to make detailed comparisons between the present results and those on the triesters, since the absolute configurations of the modified phosphate groups in the triesters are not known. FIG. 13d shows the CD spectrum of a 1:1 mixture of the diastereoisomers of d-TpT. The spectrum of this mixture is clearly different than the spectrum of d-TpT (Cantor et al., 1970). For d-TpT positive (θ) occurs at 280 nm with a magnitude approximately 1.8-fold greater than the 275 nm band of d-TpT. Similarly, d-TpT shows negative (θ) at 250 nm which is approximately 1.8-fold greater than the negative band at 245 nm in d-TpT. Interaction of Dideoxyadenosine Methyl Phosphonates with Poly (U) and Poly (dT) Both diastereoisomers of d-ApA form complexes with poly U at 0° C. The mixing curves for d-ApA 1 and d-ApA) 2 with poly U (FIG. 14) show that complex formation occurs with a base stoichiometry of 2U:1A. Similar results were obtained for the interaction of d-ApAp) 1 and d-ApAp) 2 with poly U and for the interaction of d-ApA with poly (dT). As shown in FIG. 15, the methyl phosphonate-polynucleotide complexes exhibit a cooperative thermal transition with a well-defined melting temperature. The melting temperature of the d-ApA) 1 -poly U complex is 4.4° higher than the d-ApA) 2 -poly U complex. A similar difference in melting temperatures for the d-ApAp-poly U complexes was also observed (Table 5). Essentially no difference is observed between the Tm values of the two d-ApA-poly(dT) complexes, however. Significant increases are observed in the thermal stabilities of the dinucleoside methyl phosphonate-polynucleotide complexes as compared to similar complexes formed between d-ApA and poly U or poly dT. The non-ionic d-Apa forms complexes with Tm values 8.4° and 12.4° higher that that of d-Apa-poly U, while the singly-charged d-ApAp from complexes with Tm values 6.5° and 10.4° higher than d-Apa-poly U. Similarly, the complexes formed between d-ApA and poly (dT) each melt approximately 10° higher than the d-ApA-poly(dt) complex. TABLE 5______________________________________Melting Temperatures of Complexes Formed BetweenDideoxyadenosine Methyl Phosphonate Analogs andPolyuridylic Acid or Polythymidylic AcidComplex.sup.(a) Tm °C. (poly U).sup.(b) Tm °C. (poly dT)______________________________________d-ApA 7.0 9.2(d-ApA).sub.1 15.4 18.7(d-ApA).sub.2 19.8 18.4(d-ApAp).sub.1 13.5 --(d-ApAp).sub.2 17.4 --______________________________________ .sup.(a) Complex stoichiometry: 2U:1A or 2T:1A .sup.(b) 10 mM Tris.HCl 10 mM MgCl.sub.2 pH 7.5 Discussion Dinucleoside methyl phosphonates are novel nucleic acid analogs in which the phosphodiester internucleoside linkage is replaced by a 3'-5' linked internucleoside methyl phosphonyl group. Unlike the dinucleoside methylene phosphonates prepared by others, the methyl phosphonate analogs do not contain a negatively charged backbone and are nonionic molecules at pH 7. The methyl phosphonate group is isosteric with respect to the phosphate group of dinucleoside monophosphates. Thus, these analogs should present minimal steric restrictions to interaction with complementary polynucleotides or single-stranded regions of nucleic acid molecules. Since the methyl phosphonyl group is not found in naturally occuring nucleic acid molecules, this internucleoside linkage may be resistant to hydrolysis by various nuclease and esterase activities and this has in fact been observed (Miller, unpublished data). These properties make analogs of this type potentially useful as vehicles for exploring the interactions of selected oligonucleotide sequences with nucleic acids and nucleic acid-related enzymes within the living cell (Miller et al., 1977). The preparation of the oligonucleoside methyl phosphonates follows the basic strategy used for the preparation of protected oligonucleotide phosphotriesters. The synthetic scheme which has been adopted first involves preparation of a protected nucleoside 3'-methyl phosphonate β-cyanoethyl ester (FIG. 17). This two-step preparation can be carried out in a one-flask reaction and proceeds in high overall yield. Since the product is readily purified by silica gel column chromatography, multigram quantities of this key intermediate can be prepared. By selective removal of the 5'-trityl group or the β-cyanoethyl group, chain extension can proceed in either direction. Thus, compound 2 in FIG. 17 serves as a basic building block for the preparation of longer oligomers. This type of synthetic scheme was originally developed by others for the preparation of oligonucloetide βββ-trichloroethyl phosphotriesters and has more recently been used by others for the preparation of oligonucleotide p-chlorophenyl phosphotriesters. This procedure also allows the preparation of specifically [ 13 C]enriched dimers by use of [ 13 C]-methyl phosphonic acid in the synthesis of 1. Dimers and oligomers containing [ 13 C]-methyl phosphonate groups could be very useful for probing the physical and biological properties of oligonucleoside methyl phosphonates by nuclear magnetic resonance spectroscopic techniques (Cheng et al., manuscript in preparation). In the present study (FIG. 17), the -cyanoethyl group was removed from 1 and chain extension was continued in the 3'-direction. Two types of condensation reactions were carried out: (1) condensation with a 3'-O-acetylated nucleoside to give dimers with the general structure 6 and (2) condensation with a -ucleoside 3'-methyl phosphonate β-cyanoethyl ester to give dimers with general structure 7. The latter type of dimer can be further extended by removal of the β-cyanoethyl group followed by condensation with other oligonucleoside methyl phosphonate blocks. In this way, oligonucleoside methyl phosphonates containing up to four deoxyadenosine residues and up to nine thymidine residues have been prepared. Different condensing agents were used in these reactions, including dicyclohexylcarbodiimide (DCC), triisopropylbenzenesulfonyl chloride (TPSCl) and mesitylenesulfonyl tetrazolide (MST). The order of condensing efficiency was found to be MST>TPSCl>DCC. Although DCC did bring about condensation, several days at elevated temperatures were required and the yields were quite low. Considerable improvement in reaction yield was obtained when TPSCl was used. However, again prolonged reaction periods were required and noticeable buildup of side products was observed. The reagent of choice for these reactions in MST. The reaction occurs within a period of several hours, with little or no side products. The efficiency of a particular condensing agent depends not only upon its structure but also upon the nature of the phosphorous-containing substituent which is activated. Thus, when MST was used as a condensing agent, we observed that reactions involving nucleoside 3'-methyl phosphonates or nucleoside 3'-ethyl phosphates usually proceed in lower yield than those involving nucleoside 3'-p-chlorophenyl phosphates. The ability to separate the individual diastereoisomers of each dimer sequence allowed examination of effect of the configuration of the phosphonyl methyl group on the overall dimer conformation. As shown in FIG. 16, the isomers differ in configuration at the internucleoside linkage with the methyl group assuming either a pseudo-axial or pseudoequitorial position when the dimers are drawn in a stacked conformation. The unique conformational properties of each diastereoisomer of d-ApA and d-ApAp are most readily seen by examining the percent hypochromicity of each diastereoisomer (Table 4). Isomer 1 of both d-ApA and d-ApAp exhibits a greater percent hypochromicity than does isomer 2 of this series. Since the percent hypochromicity is related to the extent of base-base overlap in dimers of this type (Ts'o, 1974), the result suggests that d-ApA) 1 and d-ApAp) 1 and more highly stacked in solution than are d-ApA) 2 and d-ApAp) 2 . Comparison of the percent hypochromicities of the methyl phosphonate dimers with that of d-ApA shows that these dimers are less stacked than the parent dinucleoside monophosphate. A similar result was observed for the methyl and ethyl phosphotriesters of d-ApA. Thus, non-ionic methyl or ethyl phosphotriester or methyl phosphonate internucleoside linkages appear to perturb the stacking interactions between the bases in these dimers. The circular dichroism spectra of dinucleoside monophosphates are indicators of both the extent and mode of base stacking, as well as the population of right-handed versus left-handed stacks. The CD spectra of each diastereoisomer for the methyl phosphonate dimer sequences d-ApA, d-ApAp, d-ApT and d-TpA suggest that each diastereoisomer has an unique stacking mode in solution. The profiles of the CD spectra of d-ApA) 1 and d-ApAp) 1 are very similar to those of d-Apa and r-ApA, and differ only in the magnitude of the molecular ellipticity. This result and the results of the hypochromicity measurements suggest that the stacking modes of the bases in these dimers are similar to those of d-ApA and r-ApA. On the other hand, the profiles of the DC spectra of d-ApA) 2 and d-ApAp) 2 are quite different. The magnitudes of the molecular ellipticities of dApA) 2 and dApAp) 2 are greatly diminished, with complete loss of [θ] at 270 nm. Since the hypochromicity measurements suggest that the bases in these dimers have substantial overlap, the mode of stacking in these dimers must be quite different from that found for isomer 1 or for d-ApA. The magnitude of the molecular ellipticity in dimers of this type is sensitive to the angle, θ, between the transition dipoles of the bases. The value of the molecular ellipticity is greatest when θ is 45° and diminishes to 0 when θ is 0°, 90° or 180°. Thus, the most reasonable interpretation of the CD results is that in d-ApA) 1 and d-ApAp) 1 , the bases tend to orient in an oblique manner, while in d-ApA) 2 and d-ApAp) 2 , the bases tend to orient in a parallel or perpendicular manner. This interpretation is supported by the base-base stacking patterns as determined by pmr spectroscopy. The substantial change in the CD profile of d-ApA) 2 rather than a simple diminution of the amplitude of the [θ] values suggests that variation of the population of right-handed versus left-handed stacks would not provide an adequate explanation of the CD results. The CD spectra of d-ApT isomers 1 and 2 have the same shape as the CD spectrum of d-ApT, but with diminished molecular ellipticity. For d-TpA, the spectrum of isomer 2 is identical to that of d-TpAp, while the spectrum of isomer 1 shows diminished [θ] values of the peak and trough regions. Thus, the stacking modes in these methyl phosphonate dimers are expected to be basically similar to the stacking modes of the parent dinucleoside monophosphates, but with perhaps different degrees of base-base overlap or different populations of right- and left-handed stacks. The dimer, d-ApA, forms stable complexes with both polyribo- and polydeoxyribonucleotides. These poly U and poly dT complexes have greater stability than similar complexes formed by the parent dinucleoside monophosphate, d-ApA. Similar observations have previously been made for triple helix formation between the alkyl phosphotriesters d-Ap(Me)A or d-Ap(Et)A and poly U, for duplex formation between oligonucleotide triesters and tRNA and for helical duplex formation between the octathymidylate ethyl phosphotriester, d-[Tp(Et)] 7 T, and poly dA. It should be noted, however, that d-[Tp(Et)] 7 T, in contrast to d-ApA, exhibits selective binding to polydeoxyribonucleotides versus polyribonucleotides in duplex formation. Previous analyses indicate that the increased stability of the complexes formed between nonionic oligomers and complementary polynucleotides results from the reduction in charge repulsion between the nonionic backbone of the oligomer and the negatively charged sugar-phosphate backbone of the polynucleotide. Although both d-ApAp and d-ApA possess a formal negative charge, the d-ApAp.poly U complexes are more stable than the d-ApA.poly U complex. The 3'-terminal methyl phosphonate group of dApAp is free to rotate away from the negatively charged phosphate backbone of poly U without disrupting the base-pairing and base-stacking interactions in the complex. In contrast, repulsion between the negative charge of the phosphodiester linkage in d-ApA and the polymer backbone directly opposes base-pairing and stacking. Thus, the presence of a negative charge at the internucleotide linkage contributes much more effectively to the charge repulsion effect between the dimers and polynucleotides. Under the conditions of the present experiments, the Tm values of d-Ap(Me)A.poly U and d-Ap(Et)A.poly U are 13° C. and 12° C. respectively. These Tm values are lower than those of d-ApA and d-ApAp complexes with poly U. These results suggest that the increasing size of the methyl and ethyl side chains in the phosphotriester dimers may provide a greater steric hindrance to complex formation. The methyl group of the phosphonate dimers should be only slightly larger in size than the oxygen of the phosphate group, and thus would be expected to have the least steric effect. A similar phenonenon has been observed when the stabilities of poly U complexes with the ethyl phosphotriester and methyl phosphonate analogs of d-ApApApA are compared. The differences in the conformations of the individual diastereoisomers of d-ApA and d-ApAp are reflected in their interactions with poly U. For each dimer, the diastereoisomer with greater base-base overlap (isomer 1) forms a complex of lower stability with poly U. In a previous analysis of the influence of C-2' substituents of adenine polynucleotides on the Tm values of the helices, it can be reasoned that the conformation free-energy difference (F D -F S ) at the melting temperature is directly related to the Tm value, where F D represents the free energy of the double-stranded duplex, and F S represents the free energy of the base-stacked single strand. The values of F D -F S reflect the conformation of the duplex state and the single-stranded state. The data indicates that (F D -F S ) for isomer 1 of dApA or dApAp is slighty less than (F D -F S ) for isomer 2 of dApA or dApAp. This reduction may reflect a higher F S value of isomer 1 since this isomer indeed has a greater degree of stacking, assuming that F.sub. D for isomer 1 and isomer 2 remains the same. In contrast to the behavior with poly U, both diastereoisomers of d-ApA form complexes with poly(dT) which have similar Tm values. Since the geometry of the triple helix of dApA.2 poly U is likely to be different than the geometry of the dApA.2 poly dT triple helix, the difference in F S of isomer 1 versus F S of isomer 2 may be compensated by a difference in F D of isomer 1 versus F D of isomer 2. The studies reported here have shown that dideoxyribonucleotide analogs containing nonionic 3'-5' internucleoside methyl phosphonate linkages can be readily synthesized. The configuration of the methyl group in the backbone of these dimers influences their conformation in solution and their ability to form complexes with complementary polyribonucleotides. In addition, preliminary studies have shown that oligodeoxyribonucleoside methyl phosphonates are resistant to nuclease hydrolysis, are taken up in intack form by mammalian cells in culture and can exert specific inhibitory effects on cellular DNA and protein synthesis. Unlike 2'-O-methyl oligonucleotide ethyl phosphotriesters, the methyl phosphonates appear to have relatively long half-lives within the cells. Thus, oligonucleoside methyl phosphonates of specific sequence could complement oligonucleotide phosphotriesters as probes and regulators of nucleic acid function within living cells. Nonionic Nucleic Acid Alkyl and Aryl Methylphosphonates There will now be described the synthesis of a series of oligonucleoside methylphosphonates whose base sequences are complementary to the anticodon loops of tRNA lys and to the --ACCAOH amino acid accepting stem of tRNA. The effects of these analogues on cell-free aminoacylation and cell-free protein synthesis will be considered. The uptake of selected analogues by mammalian cells in culture and the effects of these compounds in bacterial and mammalian cell growth are also discussed. Materials Nucleosides were checked for purity by paper chromatography before use. N-Benzoyldeoxyadenosine, N-isobutyryldeoxyguanosine, their 5'-O-dimethoxytrityl derivatives and 5'-O-monomethoxytrityl thymidine were prepared according to published procedures. d-[(Meo) 2 Tr]bzApbzApCE, d-[(Meo) 2 Tr]bzApbzAOAC, d-[(Meo)Tr]TpTpCE, d-ApT, d-Ap-[ 3 H]-T, d-TpT and d-Tp-[ 3 H]T were also synthesized by standard procedures. Dimethylmethylphosphonate and benzenesulfonic acid were used without further purification. Hydracrylonitrite was dried over 4 Å molecular sieves. Methylphosphonic acid dipyridinium salt and mesitylenesulfonyl tetrazolide were prepared. Anhydrous pyridine was prepared by refluxing reagent grade pyridine (3L) with chlorosulfonic acid (40 ml) for 7 hrs followed by distillation onto sodium hydroxide pellets (40 g). After refluxing for 7 hrs, the pyridine was distilled onto 4 Å molecular sieves and stored in the dark. Silica gel column chromatography was carried out using Baker 3405 silica gel (60-200 mesh). Thin layer silica gel chromatography (TLC) was performed on E. Merck Silica Gel 60F 254 plastic backed TLC sheets (0.2 mm thick). High pressure liquid chromatography (HPLC) was carried out using a Laboratory Data Control instrument on columns (2.1 mm×1 m) packed with HC Pellosil. The columns were eluted with a linear gradient (40 ml total) of chloroform to 20% (V/V) methanol in chloroform at a flow rate of 1 ml/min. Ultraviolet spectra were recorded on a Cary 14 or a Varian 219 ultraviolet spectrophotometer with a thermostatted cell compartment. The following extinction coefficients (260 nm) were used: d-T, 9,100; d-[(Meo)Tr]T, 10,200; d-[(Meo) 2 Tr]bzA,12,500; d-bzA, 10,600; d-[(Meo) 2 Tr]ibuG, 17,400; and d-ibuG, 16,700. Paper chromatography was carried out on Whatman 3 mm paper using solvent A: 2-propanol-concentrated ammonium hydroxidewater (7:1:2 V/V). Preparation of d-[(Meo) 2 Tr]ibuGpCE: d-[(Meo) 2 Tr]ibuG (12 g; 18.7 mmoles) and the pyridinium salt of methyl-phosphonic acid (21 mmoles) were dried by evaporation with anhydrous pyridine (4×20 ml) and the residue in 40 ml of pyridine was treated with 2,4,6-triisopropylbenzenesulfonyl chloride (12.7 g, 42 mmoles) for 8 hrs at room temperature. Hydracrylonitrile (4.5 g, 63 mmoles) and 2,4,6-triisopropylbenzenesulfonyl chloride (0.61 g, 2 mmoles) were added and the reaction mixture was kept at room temperature. After 2 days the reaction mixture was poured into 500 ml of ice-cold 5% NaHCO 3 solution. The solution was extracted with ethyl acetate (2×250 ml) and the combined extracts were dried over anhydrous Na 2 SO 4 . Examination of the extract by TLC showed the presence of both d-[(Meo) 2 Tr]ibuGpCE (Rf-0.31 silica gel tlc, 10% MeOH/CHCl 3 ) and d-ibuGpCE (Rf-0.14, silica gel tlc, 10% MeOH/CHCl 3 ). After concentration, the ethyl acetate extract was chromatographed on silica gel (4×35 cm) using ether (1L ) and a 0 to 20% linear gradient of methanol in chloroform (1.6L total) as solvents. d-[(Meo) 2 Tr]ibuGpCE (2.75 mmoles) was obtained in 15% yield while d-ibuGpCE (2.46 mmoles) was obtained in 13% yield. Additional d-[(Meo) 2 Tr]ibuGp (3.69 mmoles, 20%) was obtained from the aqueous bicarbonate solution after extraction with chloroform (2×200 ml). Preparation of Protected Oligonucleoside Methylphosphonates: The same general procedures were used for the preparation of dinucleoside methylphosphonates. The specific conditions used in the condensation reactions and the yields obtained after silica gel column chromatography are given in Table VI. The ultraviolet spectroscopic characteristics and the mobilities of the protected oligonucleotides on silica gel TLC and silica gel HPLC are given in Table VII. Preparation of Oligonucleoside Methylphosphonates: The protecting groups were removed from the blocked oligonucleoside methylphosphonates using conditions described previously. In the case of the dA-containing oligomers, the N-benzoyl groups were removed by treatment with hydrazine. The oligomers were purified by preparative paper chromatography using solvent A. For the [ 3 H] -labeled oligothymidine methylphosphonates, d-Tp) n -[ 3 H]-T, the condensation reactions containing d-[(Meo) 2 Tr]Tp) n +[ 3 H]TOAC were run on 0.01 (n=1) and 0.005 (n=4,8) mmole scales while d-GpGp-[ 3 H]-T was prepared on a 0.012 mmole scale. After completion of the reaction, the protecting groups were removed and the entire reaction mixture was chromatographed on paper. The oligonucleoside methylphosphonates were eluted from the paper with 50% aqueous ethanol. The ethanol solutions were passed through DEAE cellulose columns (0.5×1 cm) and stored at 0° C. The UV spectral properties and chromatographic mobilities of the oligonucleoside methylphosphonates are given in Table VII. For use in the physical, biochemical, and biological experiments described below, aliquots containing the required amount of oligomer were evaporated to dryness and the oligomer was dissolved in the buffer used in the particular experiment. TABLE VI__________________________________________________________________________Preparation of Protected Oligodeoxyribonucleoside Methylphosphonates3'-Methylphosphonate 5'-OH MST Product YieldComponents (mmoles) Component (mmoles) (mmoles) (mmoles) %__________________________________________________________________________d-[(MeO).sub.2 Tr]ibuGp (0.50) d-ibuGpCE (0.50) 2.0 d-[(MeO).sub.2 Tr]ibuGpibuGpCE (.082) 16d-[(MeO).sub.2 Tr]ibuGp (1.0) d-bzAOAC (1.5) 4.0 d-[(MeO).sub.2 Tr]ibuGpbzAOAC 42.42)d-[(MeO)Tr]TpTp (0.33) d-TpTpCE (0.50) 1.6 d-[(MeO)Tr]TpTpTpTpCE 50.168)d-[(MeO)Tr]Tp(Tp).sub.2 TpCE (0.0324) d-Tp(Tp).sub.2 TpCE (.0524) 0.16 d-[(MeO)Tr]Tp(Tp).sub.6 TpCE 430138)d-[(MeO).sub.2 Tr]ibuGpibuGp (.07) d-TOAC (.15) 0.28 d-[(MeO).sub.2 Tr]ibuGpibuGpTOAC (0.0153) 22d-[(MeO).sub.2 Tr]bzApbzAp (0.065) d-bzAOAC (0.043) 0.163 d-[(MeO).sub.2 Tr]bzApbzApbzAOAC (0.023) 53d-[(MeO).sub.2 Tr]bzApbzAp (0.13) d-bzApbzAOAC (0.20) 0.52 d-[(MeO).sub.2 Tr]bzApbzApbzApbzAOAC (0.031) 24d-[(MeO).sub.2 Tr]bzApbzAp (0.0168) d-ibuGpbzAOAC (0.0168) 0.0735 d-[(MeO).sub.2 Tr]ApbzApibuGpbzAOAC (0.0029) 1__________________________________________________________________________ TABLE VII__________________________________________________________________________Ultraviolet Spectral Properties and Chromatographic Mobilities ofProtectedOligodeoxyribonucleoside Methylphosphonates Silica UV Spectra.sup.a Gel HPLC.sup.c λ max. λ min. ε.sub.260/235 ε.sub.260/280 Silica Gel TLC Retention timeOligomer nm nm calcd. obsvd. calcd. obsvd. 5% 10% 15% 20% (min)__________________________________________________________________________d-[(MeO)Tr]TpTpTpTpCE 265 243 1.34 1.31 1.55 1.64 -- -- .08 .29 -- 235 shd-[(MeO)Tr]Tp(Tp).sub.6 TpCE 265 243 1.75 0.92 1.57 1.56 -- 0.00 -- .13 --d-[(MeO).sub.2 Tr]ibuGpibuGpCE 238 225 1.19 1.05 1.33 1.32 -- 0.16 -- -- 19.2 253 245 260 256 280 270d-[(MeO).sub.2 Tr]ibuGpbzAOAC 235 256 0.82 0.75 0.88 0.87 -- 0.29 -- -- 12.3 278d-ibuGpbzAOAC 260 239 1.63 1.27 0.90 0.90 -- 0.18 -- -- 15.5 280 267 0.14 17.6d-[(MeO).sub.2 Tr]ibuGpibuGpTOAC 240 sh 228 1.34 1.51 1.38 1.45 -- 0.18 -- -- 16.0 260 275 shd-[(MeO).sub.2 Tr]bzApbzApbzAOAC 234 227 0.66 0.61 0.59 0.59 -- 0.41 0.55 -- 13.4 280 255 0.38 0.53 14.3d-[(MeO).sub.2 Tr]bzApbzApbzApbzAOAC 233 sh 253 0.71 0.60 0.59 0.60 -- -- 0.31 -- 19.3 280d-[(MeO).sub.2 Tr]bzApbzApibuGpbzAOAC 235 sh 255 0.89 0.74 0.73 0.75 -- 0.15 0.44 -- 23.8 280__________________________________________________________________________ .sup.a Measured in 95% ETOH .sup.b EM silica gel 60 F.sub.254 sheets, 0.2 mm thick. .sup.c HC Pellosil (2.1 mm × 1 m) 0% to 20% methanol in chloroform ml/min., 40 ml total volume Interaction of Oligodeoxyadenylate Methylphosphonates With Polynucleotides The continuous variation experiments and melting experiments were carried out. The extinction coefficients of the oligomers were determined by comparing the absorption of a solution of the oligomer in water at pH 7.0 to the absorption of the same solution at pH 1.0. The oligomer extinction coefficient was calculated from the observed hyperchromicity of the oligomer at pH 1.0 by using the following extinction coefficients: d-A pH 1.0, 14.1×10 3 and d-G pH 1.0, 12.3×10 3 . The molar extinction coefficient of poly(U) is 9.2×10 3 (265 nm) and of poly(dT) is 8.52×10 3 (264 nm). Cell-Free Aminoacylation (1) E. coli system: Unfractionated tRNA E. coli was purchased from Schwarz Mann and unfractionated E. coli aminoacyl synthetase was purchased from Miles Laboratories, Inc. Reactions were run in 60 μl buffer containing 100 nM Tris, HCl, pH 7.4, 10 mM Mg(OAC) 2 , 5 mM KCl, 2 mM ATP, 4 μM [ 3 H]-amino acid, 1.8 μM tRNA coli and 0 to 100 μM oligonucleotide. Reactions were initiated by addition of 4 μg of aminoacyl synthetase. Aliquots (10 μl) were removed at various times, added to 1 ml of cold 10% trichloroacetic acid and the resulting precipitate filtered on Whatman G/F filters. After washing with 4 (1 ml) portions of 2N HCl and 4 (1 ml) portions of 95% ETOH, the filters were dried and counted in 7 ml New England Nuclear 949 scintillation mixture. (2) Rabbit Reticulocyte System A rabbit reticulocyte cell-free translation system was obtained from New England Nuclear. Reactions were run in 12.5 μl of buffer containing 1 μl translation mixture, 79 mM potassium acetate, 0.6 mM magnesium acetate, 57 μM [ 3 H]-Lysine, and 50 μM oligomer. The reactions were initiated by addition of 5 μl of reticulocyte lysate and were assayed as described for the E. coli system. Cell-Free Protein Synthesis (1) E. coli system A cell-free protein synthesizing system was isolated from E. coli B cells (S-30). The system incorporates 300 pmoles of [ 3 H]-phenylalanine/mg of S-30 protein after 15 min incubation at 37° C. when poly U is used as a message. (2) Rabbit Reticulocyte The reticulocyte translation system prepared by New England Nuclear was used. For the translation of globin mRNA the reactions were run in 12.5 μl of buffer containing: 1 μl of translation mixture, 0.10 μg of globin mRNA (Miles Laboratories), 79 mM, potassium acetate, 0.2 mM magnesium acetate 0 to 50M oligomer and 20.5 μM [ 3 H]-leucine. For the translation of poly(U) the reactions were run in 12.5 μl buffer containing: 1 μl of translation mixture, 120 mM potassium acetate, 0.8 mM magnesium acetate, 367 μM poly(U), 0 to 200 μM oligomer (base concentration) and 32 μM [ 3 H]-phenylalanine. Reactions were initiated by addition of 5 μl of reticulocyte lysate. Aliquots (2 μl) were removed at various times and added to 1.0 ml of bovine serum albumin (100 μg) solution. The protein was precipitated by heating with 1 ml of 10% trichloroacetic acid at 70° C.; filtered on G/F filters and counted in 7 ml of Betaflour. Uptake of Oligodeoxyribonucleoside Methylphosphonates The uptake of d-Ap-[ 3 H]-T, d-GpGp-[ 3 H]-T and d-Tp) n -[ 3 H]-T by transformed Syrian hamster fibroblasts were determined. Effects of Oligodeoxyribonucleoside Methylphosphonates On Colony Formation (1) E. coli E. coli B was grown in M-9 medium supplemented with glucose (36 g/l) and 1% Casamino acids. The cells were harvested in midlog phase and resuspended in 50 μl of fresh medium containing 0 to 160 μM oligomer at a final cell density of 1×10 4 cells/ml. The cells were incubated for 1 hr. at 37° C. and then diluted with 0.9 ml of medium. A 0.8 ml aliquot was added to 2.5 ml of 0.8% Bactoagar at 45° C. This solution was quickly poured onto a 100 mm plate containing solid 1.2% Bactoagar. After solidification, the plates were incubated overnight at 37° C. and the resulting colonies were counted. (2) Transformed Syrian Hamster Embryonic Fibroblasts (BP-6) and Transformed Human Fibroblasts (HTB1080) Colony formation by the fibroblasts in the presence of the methylphosphonate analogues was carried out. RESULTS Synthesis of Oligodeoxyribonucleoside Methylphosphonates The synthetic scheme used for preparing the oligonucleoside methylphosphonates followed the basic approach used to synthesize dideoxyribonucleoside methylphosphonates. Suitably protected monomers or oligomer blocks carrying a 3'-terminal methylphosphonate group were condensed with protected mono- or oligonucleotides bearing a free 5'-hyroxyl group. Mesitylenesulfonyl tetrazolide was used as the condensing agent. The fully protected oligomers were purified by silica gel column chromatography. The reaction conditions used and the yields obtained are given in Table VI. The oligomers were characterized by ultraviolet spectroscopy, thin layer chromatography and high pressure liquid chromatography as indicated in Table VII. The protecting groups were removed as previously described. In the case of the deoxyadenosine-containing oligomers, the N-benzoyl groups were first removed by treatment with hydrazine hydrate. The remaining 3'-O acetyl and 5'-O dimethoxytrityl groups were removed by sequential treatment with ammonium hydroxide and 80% acetic acid. The oligomers were purified by preparative paper chromatography and were characterized by UV spectroscopy (Table VIII). Interaction of Oligodeoxyribonucleoside Methylphosphonates With Complementary Polynucleotides Table IX summarizes the melting temperatures of complexes formed between oligodeoxyadenosine methylphosphonates and poly (U) or poly (dT). For comparison, the melting temperatures of complexes formed by oligodeoxyribo- and oligoriboadenosines are included. Each oligomer forms a triple-stranded complex with a stoichoimetry of 2U:1A or 2T:1A. The melting temperatures increase as the chainlength of the oligonucleotide increases. For a given chain length, the complexes formed by the methylphosphonate analogues melt at higher temperatures than those formed by the natural diester oligomers. TABLE VIII__________________________________________________________________________Spectral Properties and Chromatographic Mobilities of OligodeoxyribonucleosideMethylphosphonates UV Spectra.sup.a Paper Chromatography.sup.b λ max. λ min. ε RfOligomer nm nm ε.sub.260 /ε.sub.280 λ max. Solvent A__________________________________________________________________________d-GpGpT.sup.c 257 230 1.45 33.4 × 10.sup.3 0.31 270 shd-ApApA 258 232 4.27 39.0 × 10.sup.3 0.29d-ApApApA 258 230 3.77 50.4 × 10.sup.3 0.11d-ApApGpA 258 227 3.03 50.3 × 10.sup.3 0.11d-Tp-[.sup.3 H]--T 267 235 1.53 -- 0.59(d-Tp).sub.4 -[.sup.3 H]--T 266 235 1.49 -- 0.21(d-Tp).sub.8 -[.sup.3 H]--T 266 235 1.56 -- 0.17__________________________________________________________________________ .sup.a Measured in water, pH 7.0 .sup.b Rf.sup.A pT = 0.11 .sup.c The UV spectrum is similar to that of dGpGpT (Miller et. al., 1974). TABLE IX______________________________________Interaction of Oligonucleoside Methylphosphonates withComplementary Polynucleotides.sup.aOligomer Tm Poly U (2U:1A) Tm Poly dT (2T:1A)______________________________________d-ApA isomer 1 15.4 18.7.sup. isomer 2 19.8 18.4d-ApApA 33.0 36.8d-ApApApA 43.0 44.5d-ApA 7.0 9.2d-ApApApA 32.0 35.5r-ApApApA 36.2° C. 2.4° C.______________________________________ .sup.a 5 × 10.sup.-5 M total [nucleotide], 10 mM Tris, 10 mM MgCl.sub.2, pH 7.5 With the exception of r-ApApApA, the complexes formed by the oligomers with poly (dT) have slightly higher melting temperatures than the corresponding complexes formed with poly (U). The interaction of d-GpGp-[ 3 H]-T with unfractionated tRNA E.coli was measured by equilibrium dialysis. The apparent association constants at 0°, 22°, and 37° C. are 1,100M -1 , 200M -1 , and 100M -1 respectively. These binding constants are much lower than those of the 2'-O-methylribooligonucleotide ethyl phosphotriester, G m p (Et)-G m p(Et)-[ 3 H]-U, which are: 9,300M -1 (0° C.), 1,900M -1 , (22° C.) and 2,000M -1 (37° C.). Effect of Oligodeoxyribonucleoside Methylphosphonates on Cell-Free Aminoacylation to tRNA The effects of selected oligodeoxyribonucleoside methylphosphonates on aminoacylation of unfractionated tRNA E. coli are shown in Table X. Three amino acids were tested at various temperatures. The deoxyadenosine-containing analogs which are complementary to the --UUUU-- sequence of the anticodon of tRNA E .coli lys have the largest inhibitory effect on aminoacylation of tRNA E .coli lys . The percent inhibition increases with increasing chain length and decreases with increasing temperature. Inhibition by d-ApApGpA and by the diesters dApApApA is less than that exhibited by d-ApApApA. In contrast to their behavior with TRNA E .coli lys , neither the methylphosphonates, d-ApApApA, d-ApApGpA nor the phosphodiesters, d-ApApApA, r-ApApApA, had any inhibitory effect on tRNA rabbit lys in the rabbit reticulocyte cell-free system (data not shown). TABLE X______________________________________Effects of Oligonucleoside Methylphosphonates onAminoacylation in an E. coli Cell-Free System % Inhibition.sup.(b) Phe Leu LysineOligomer.sup.(a) 0° C. 0° C. 0° C. 22° C. 37° C.______________________________________d-ApA 6 0 7 -- --d-ApApA 9 0 62 15 0d-ApApApA 9 12 88 40 16d-ApApGpA 12 12 35 0 --d-GpGpT 31 5 34 9 15dGpGpT (400 μM) 23 -- -- -- --d-ApApApA 0 7 .sup. 71.sup.(c) .sup. --5.sup.(c)r-ApApApA -- -- .sup. 78.sup.(d) .sup. --7.sup.(d)______________________________________ .sup.(a) [oligomer] = 50 μM .sup.(b) [tRNA.sub.coli ] = 2 μM .sup.(c) [oligomer] = 100 μM .sup.(d) [oligomer] = 125 μM Effects of Oligodeoxyribonucleoside Methylphosphonates On Cell-Free Protein Synthesis The ability of deoxyadenosine containing oligonucleoside methylphosphonates to inhibit polypeptide synthesis in cell-free systems directed by synthetic and natural messages was tested. The results of these experiments are given in Table XI. Poly (U) directed phenylalanine incorporation and poly (A) directed lysine incorporation are both inhibited by oligodeoxyadenosine methylphosphonates and diesters in the E.coli system at 22° C. The percent inhibition increases with oligomer chain length and is greater for polyphenylalanine synthesis. The methylphosphonate analogues are more effective inhibitors than either d-ApApApA or r-ApApApA at the same concentration. Although both the oligodeoxyadenosine methylphosphonates and the phosphodiesters inhibit translation of poly (U) in the rabbit reticulocyte system, no effect on the translation of globin message was observed. As in the case of the E.coli system, inhibition of phenylalanine incorporation increased with oligomer chain length and was greater for the methylphosphonate analogues than for the diesters. Uptake of Oligodeoxyribonucleoside Methylphosphonates By Mammalian Cells FIG. 18 shows the incorporation of radioactive 100 μM d-GpGp-[ 3 H]-T with time into transformed Syrian hamster embryonic fibroblasts growing in monolayer. The incorporation is approximately linear for the first hour and begins to level off after 1.5 hours. The concentration of radioactivity inside the cell is approximately 117 μM after 1.5 hours assuming a cell volume of 1.5 μl/10 6 cells. TABLE XI______________________________________Effects of Oligonucleoside Methylphosphonates on Bacterialand Mammalian Cell-Free Protein Synthesis at 22° C. E. coli Rabbit Reticulocyte Poly U Poly A Poly U Globin mRNAOligomer directed.sup.(a) directed.sup.(b) directed.sup.(a) directed.sup.(c)______________________________________d-ApA 20 10 -- --d-ApApA 84 30 81 --d-ApApApA 100 65 77 0d-ApApGpA 22 -- -- 0d-ApApApA 13 19 18 0r-ApApApA 18 17 85 0______________________________________ .sup.(a) [Poly U] = 360 μM in U [oligomer] = 175-200 μM in base .sup.(b) [Poly A] = 300 μM in A [oligomer] = 175-200 μM in base .sup.(c) [oligomer] = 200 μM in base Cells were incubated with 25 μM d-GpGp-[ 3 H]-T for 18 hours. The medium was removed, the cells were washed with phosphate buffer and then lysed with SDS. Approximately 30% of the total radioactivity from the lysate was found in TCA precipitable material. The DNA was precipitated from the lysate and digested with deoxyribonuclease and snake venom phosphodiesterade. The culture medium, the DNA-free lysate and the DNA digest were each examined by paper chromatography. Only intact d-GpGp-[ 3 H]-T was found in the medium. Radioactivity corresponding to [ 3 H]-TTP (6%) and to d-GpGp-[ 3 H]-T (94%) was found in the lysate, while the DNA digest gave [ 3 H]-dpT and [ 3 H]-dT as products. Similar uptake studies were carried out with d-Ap[ 3 H]-T and with a series of oligothymidylate analogues, d-(Tp)n-[ 3 H]-T (n=1,4,8). The rates and extents of uptake of these analogues were very similar to that of d-GpGp-[ 3 H]-T (FIG. 1). Examination of the culture medium and cell lysate after overnight incubation with these oligonucleotides gave results similar to those found for d-GpGp-[ 3 H]-T. Effects of Oligodeoxyribonucleoside Methylphosphonates On Colony Formation By Bacterial and Mammalian Cells The effects of selected oligodeoxyribonucleoside methylphosphonates on colony formation of E. coli B, transformed Syrian hamster fibroblast (BP-6) and transformed human fibroblast (HTB 1080) cells are summarized in Table XII. The d-Ap)nA analogues appear to inhibit E. coli colony formation at high concentrations (160 μM). However, no inhibitory effects on cellular protein or DNA synthesis could be detected in the presence of these compounds by the present assay procedures. TABLE XII______________________________________Effects of Oligonucleoside Methylphosphonates on ColonyFormation by Bacterial and Mammalian Cells in Culture % Inhibition.sup.(a) E. coli B HTB 1080Oligomer 50 μM 160 μM BP-6 (50 μM) (50 μM)______________________________________d-ApT 4 5 5, 16.sup.(b) 12d-ApA 8 58 6, <1.sup.(b) 5d-ApApA 3 44 29 31d-ApApApA 19 78 36 19d-GpGpT 7 11 7 9______________________________________ .sup.(a) The results are the average of two or three experiments. Each experiment consisted of 2 plates (bacterial cells) or 3 plates (Mammalian cells). The average variation is: ±3% in % inhibition. The cells were treated with and grown in the presence of the oligomer at 37° C. .sup.(b) The % inhibition of isomer 1 and 2 respectively. Colony formation of both transformed hamster and human cells are inhibited to various extents by the oligonucleoside methylphosphonates. Both the hamster and human cells appear to be affected to a similar extent by a given analogue. It appears in the case of dApA, that each diastereoisomer exerts a different inhibitory effect on the growth of the hamster cells. As in the case of E. coli, no inhibition of cellular protein synthesis could be detected. DISCUSSION Oligodeoxyribonucleoside methylphosphonates with sequences complementary to the anticodon loop of tRNA lys and to the --ACCA--OH amino acid accepting stem of tRNA were prepared in a manner similar to that used to prepare dideoxyribonucleoside methylphosphonates. The present studies demonstrate the ability to join blocks of protected methyl-phosphonates to give oligomers with chain lengths up to nine nucleotidyl units. The yields in these condensation reactions are acceptable, although reactions involving deoxyguanosine residues appear to proceed in low yield. Similar difficulties have been encountered in the syntheses of oligonucleotide phosphotriesters. Unlike the dideoxyribonucleoside methyl-phosphonates previously reported, the oligodeoxyribonucleoside methylphosphonates prepared for this study were not resolved into their individual diasteroisomers. The oligodeoxyadenosine analogues form triple stranded complexes with both poly(U) and poly(dT). These complexes are more stable than similar complexes formed by either oligoribo- or oligodeoxyribonucleotides. As previously suggested for oligonucleotide ethyl phosphotriesters, and dideoxyribonucleoside methyl-phosphonates, this increased stability is attributed to the decreased charge repulsion between the nonionic backbone of the analogue and the negatively charged complementary polynucleotide backbone. With the exception of r-ApApApA (Table IX), the stability of the complexes formed with poly (dT) are slightly higher than those formed with poly(U), a situation which is also observed for the interaction of poly(dA) with poly(dT) and with poly(U). The lower stability of the (r-ApApApA).2 poly(dT) complex is also reflected at the polymer level. Thus, under the conditions of the experiments described in Table IX, it was found that the Tm of poly(rA).2 poly(rU) is 83° C. while the Tm of poly(rA).2 poly(dT) is 59° C. It was observed that formation of the poly(rA).2 poly(dT) complex occurs only at a sodium ion concentration of 2.5M in the absence of magnesium, while poly(rA).2 poly(rU) forms in 0.1M sodium phosphate buffer. The oligodeoxyadenosine methylphosphonates and their parent dieters selectively inhibit cell-free aminoacylation of tRNA E .coli lys . The extent of inhibition is temperature dependent and parallels the ability of the oligomers to bind to poly(U). These observations and the previously demonstrated interaction or r-ApApApA with tRNA E .coli lys suggest that inhibition occurs as a result of oligomer binding to the --UUUU-- anticodon loop of the tRNA. The reduced inhibition observed with d-ApApGpA is consistent with this explanation, since interaction of this oligomer with the anticodon loop would involve formation of a less stable G.U base pair. Recent studies by others have shown that the rate of aminoacylation of tRNA E .col lys substituted with 5-fluorouracil is considerably lower than that of non-substituted tRNA E .coli lys . The increased Km of the 5-fluorouracil substituted tRNA suggested a decreased interaction with the lysyl aminoacyl synthetase. These results and those of others suggest that the anticodon loop of tRNA E .coli lys is part of the synthetase recognition site. Thus, inhibition of aminoacylation by the oligodeoxyribonucleoside methylphosphonates could result from the reduction in the affinity of the synthetase for tRNA lys -oligonucleotide complexes. The greater inhibition observed with d-ApApApA versus the diesters, d-ApApApA or r-ApApApA may result from greater binding of the analogue to the anticodon loop or to the decreased ability of the synthetase to displace the nonionic oligonucleotide analogue from the anticodon loop. Alternatively, oligomer binding to the anticodon loop could induce a conformational change in the tRNA, leading to a lower rate and extent of aminoacylation. Such conformational changes have been detected when r-ApApApA binds to tRNA E .coli lys . None of the oligomers have any effect on the aminoacylation of tRNA rabbit lys in a cell free system. Since the anticodon regions of tRNAs from bacterial and mammalian sources probably are similar, the oligo A analogues are expected to interact with the anticodon region of both tRNA lys s. The failure to observe inhibition of aminoacylation of tRNA rabbit lys in the presence of these oligo d-A analogs suggestes that there may be a difference between the interaction of the lysine aminoacyl synthetase with tRNA lys from E. coli and from rabbit systems, or a difference between the structure of these two tRNA lys s in response to the binding of oligo A analogues. The trimer, dGpGpT, inhibits both phenylalanine and lysine aminoacylation at 0°, but has little effect on leucine aminoacylation. The aminoacyl stems of both tRNA E .coli lys ; and tRNA E .coli Phe terminate in a G-C base pair between nucleotides 1 and 72, while a less stable G-U base pair is found at this position in tRNA E .coli leu . Thus the observed differences in inhibition of aminoacylation by d-GpGpT may reflect differences in the ability of this oligomer to bind to the different --ACC-- ends of the various tRNAs. Inhibition of lysine aminoacylation by dGpGpT is very temperature sensitive and parallels the decrease in binding to tRNA with increasing temperature. This behavior of d-GpGpT contrasts that of G p m (ET)G p m (ET)U. Although both oligomers can potentially interact with the same sequences in tRNA, the 2'-O-methylribotrinucleotide ethyl phosphotriester binds more strongly and more effectively inhibits aminoacylation. The differences in binding ability may be due to overall differences in the conformation of the deoxyribo- versus 2'-O-methylribo backbones of these oligomers. The oligodeoxyribonucleoside methylphosphonates effectively inhibit polyphenylalanine synthesis in cell-free systems derived from both E. coli and rabbit reticulocytes. In the E. coli system, the extent of inhibition by the oligodeoxyadenosine analogures parallels the Tm values of the oligomers wiht poly(U), The tetramer, d-ApApGpA which would have to form a G.U base pair with polyU, was 4.5-fold less effective than d-ApApApA. These results suggests that the oligomers inhibit polypeptide synthesis as a consequence of forming complexes with the poly(U) message. A similar inhibitory effect by poly(dA) on the translation of poly(U) has been observed by others. It is unlikely that inhibition results from non-specific interaction of the methylphosphonates with protein components of the translation systems. In the E. coli system, poly(A) translation is inhibited to a lesser extent than is translation of poly(U), while in the reticulocyte system, no inhibition of globin mRNA translation is observed. The data suggest that the magnitude of inhibition of poly(U)-directed polypeptide synthesis in the E. coli system does not reflect proportionally the ability of the oligomer to bind to poly(U). Although the oligomer pairs d-ApApA/d-ApApApA and d-ApApApA/r-ApApApA form complexes with poly(U) which have very similar Tm's (see Table IX), in each case the methylphosphonate analogues inhibit 5.5 to 6.5 times better than do the diesters. The stronger inhibitory effect could result from a decreased ability of the ribosome to displace the nonionic oligodeoxyribonucleoside methylphosphonates form the poly(U) message, or alternatively, there may be a degradation of the oligonucleotides (phosphodiesters) by nucleases in the cellfree translation systems, but not the corresponding phosphonate analogues. Experiments with radioactively labeled oligonucleotide methylphosphonates show that these analogues are taken up by mammalian cells growing in culture. The extent of uptake is consistent with passive diffusion of the oligomer across the cell membrane. Both d-Tp-[ 3 H]-T and d-(Tp) 8 -[ 3 H]-T are taken up to approximately the same extent which suggests that there is no size restriction to uptake over this chain length range. This behavior is in contrast to results obtained with E. coli cells. Examination of lysates of mammalian cells exposed to labeled oligomers for 18 hours showed that approximately 70% of the labeled thymidine was associated with intact oligomer with the remainder found in thymidine triphosphate and in cellular DNA. These observations indicate that the oligodeoxyribonucleoside methylphosphonates, which are recovered intact from the culture medium, are slowly degraded within the cell. Failure to observe shorter oligonucleotides and the known resistance of the methylphosphonate linkage to nuclease hydrolysis suggests that degradation may result from cleavage of the 3'-terminal [ 3 H]-thymidine N-glycosyl bond with subsequent reutilization of the thymine base. The uptake process of the oligonucleoside methylphosphonates is quite different from that of previously studied oligonucleotide ethyl phosphotriesters. In the case of G p m (Et)G p m (Et)-[ 3 H]-U, the oligomer is rapidly taken up by the cells and is subsequently deethylated. Further degradation to smaller oligomers is then observed, presumably as a result of nuclease-catalyzed hydrolysis of the resulting phosphodiester linkages. Approximately 80% of the oligomer is metabolized within 24 hours. Although the rate of uptake of d-Gp(Et) Gp(ET)-[ 3 H]-T is similar to that of d-GpGp-[ 3 H]-T, examination of the cell lysate showed extensive degradation of the phosphotriester analogue. The relatively long half lives of the oligodeoxyribonucleoside methylphosphonates may be of value in potential pharmacological applications of these oligonucleotide analogues. The effects of these analogues on cell colony formation confirmed that the methylphosphonates are taken up by both mammalian and bacterial cells. All the oligomers tested inhibited colony formation of both cell types of various extents. The mechanism(s) by which these compounds exert their inhibitory effects is currently under investigation. No decrease in either overall short term cellular protein synthesis or DNA synthesis was detected by the present procedure in the presence of these compounds. This does not rule out the possibility that the syntheses of certain critical proteins are perturbed by these oligomers. Currently, studies are being made of this possibility by examination of the cellular proteins using 2-dimensional gel electrophoresis. The experiments described hereinbefore extend these studies on the use of nonionic oligonucleotides as sequence/function probles of nucleic acids both in biochemical experiments and in living cells. In a future, there will be described the effects of an oligodeoxyribonucleoside methylphosphonate complementary to the 3'-terminus of 16S rRNA on bacterial protein synthesis and growth. The results here, however, suggest that sequence specific oligonucleoside methylphosphonates may find important applications in probing and regulating nucleic acid function within living cells.

Oligonucleoside alkyl-- or aryl phosphonates are nonionic analogues of nucleic acid which possess unique physical and biological properties. These properties enable the analogues to enter living cells intact and to bind with specifically selected nucleic acids within the cell. As a result, the analogues can specifically inhibit the function or expression of a preselected nucleic acid sequence. Thus the analogues could be used to specifically inhibit the growth of tumor cells or replication of viruses in infected cells. Four methods are provided for preparing oligonucleoside methylphosphonates: (1) Coupling a protected nucleoside 3'-alkyl or aryl phosphonate with the 5'-hydroxyl group of a protected nucleoside using a condensing agent; (2) Coupling protected nucleoside 3'-alkyl or aryl phosphonic acid derivative with the 5'-hydroxy group of a protected nucleoside with the activated alkyl or aryl phosphonic acid derivative possessing functionalities which are good leaving groups; (3) Coupling a protected nucleoside 3'-alkyl or aryl phosphinate derivative with the 5'-hydroxyl group of a protected nucleoside with the resulting phosphinate derivative being then oxidized to the phosphonate; and (4) Converting a oligonucleoside methoxyphosphite derivative to the alkyl or aryl phosphonate derivative by reaction with an alkyl or aryl iodide. It has been demonstrated that procedures (1) and (2) can be used to prepare oligonucleoside methylphosphonates. Others have shown that procedure (4) can be used to prepare a diribonucleoside methylphosphonate.

CROSS REFERENCE TO RELATED APPLICATIONS [0001] The present application is a Continuation of U.S. patent application Ser. No. 13/106,640, filed May 12, 2011, now pending, which is a non-provisional of, and claims benefit of priority under 35 U.S.C. 119(e) from, U.S. Provisional Patent Application No. 61/331,103, filed May 4, 2010, and which claims priority under 35 U.S.C. 371 from PCT/IB11/01026, filed May 13, 2011, the entirety of which are expressly incorporated herein by reference. FIELD OF THE INVENTION [0002] This invention relates to the field of heatsinks or items that transfer heat between a concentrated source or sink and a fluid. BACKGROUND OF THE INVENTION [0003] A heat sink is a term for a component or assembly that transfers heat generated within a solid material to a fluid medium, such as air or a liquid. A heat sink is typically physically designed to increase the surface area in contact with the cooling fluid surrounding it, such as the air. Approach air velocity, choice of material, fin (or other protrusion) design and surface treatment are some of the design factors which influence the thermal resistance, i.e. thermal performance, of a heat sink. [0004] A heat sink transfers thermal energy from a higher temperature to a lower temperature fluid medium. The fluid medium is frequently air, but can also be water or in the case of heat exchangers, refrigerants and oil. Fourier's law of heat conduction, simplified to a one-dimensional form in the x-direction, shows that when there is a temperature gradient in a body, heat will be transferred from the higher temperature region to the lower temperature region. The rate at which heat is transferred by conduction, q k , is proportional to the product of the temperature gradient and the cross-sectional area through which heat is transferred: [0000] q k = - kA   T  x ( 1 ) [0000] where q k is the rate of conduction, k is a constant which depends on the materials that are involved, A is the surface area through which the heat must pass, and dT/dx is the rate of change of temperature with respect to distance (for simplicity, the equation is written in one dimension). Thus, according to Fourier's law (which is not the only consideration by any means), heatsinks benefit from having a large surface area exposed to the medium into which the heat is to be transferred. [0005] Consider a heat sink in a duct, where air flows through the duct, and the heat sink base is higher in temperature than the air. Assuming conservation of energy, for steady-state conditions, and applying Newton's law of cooling, gives the following set of equations. [0000] Q . = m .  c p , in  ( T air , out - T air , in ) ( 2 ) Q . = T hs - T air , av R hs ( 3 ) T air , av = T air , out + T air , in 2 ( 4 ) [0006] Using the mean air temperature is an assumption that is valid for relatively short heat sinks. When compact heat exchangers are calculated, the logarithmic mean air temperature is used {dot over (m)} the air mass flow rate in kg/s. [0007] The above equations show that when the air flow through the heat sink decreases, this results in an increase in the average air temperature. This in turn increases the heat sink base temperature. And additionally, the thermal resistance of the heat sink will also increase. The net result is a higher heat sink base temperature. The inlet air temperature relates strongly with the heat sink base temperature. Therefore, if there is no air or fluid flow around the heat sink, the energy dissipated to the air cannot be transferred to the ambient air. Therefore, the heat sink functions poorly. [0008] Other examples of situations in which a heat sink has impaired efficiency: Pin fins have a lot of surface area, but the pins are so close together that air has a hard time flowing through them; Aligning a heat sink so that the fins are not in the direction of flow; Aligning the fins horizontally for a natural convection heat sink. Whilst a heat sink is stationary and there are no centrifugal forces and artificial gravity, air that is warmer than the ambient temperature always flows upward, given essentially-still-air surroundings; this is convective cooling. [0009] The most common heat sink material is aluminum. Chemically pure aluminum is not used in the manufacture of heat sinks, but rather aluminum alloys. Aluminum alloy 1050A has one of the higher thermal conductivity values at 229 W/m·K. However, it is not recommended for machining, since it is a relatively soft material. Aluminum alloys 6061 and 6063 are the more commonly used aluminum alloys, with thermal conductivity values of 166 and 201 W/m·K, respectively. The aforementioned values are dependent on the temper of the alloy. [0010] Copper is also used since it has around twice the conductivity of aluminum, but is three times as heavy as aluminum. Copper is also around four to six times more expensive than aluminum, but this is market dependent. Aluminum has the added advantage that it is able to be extruded, while copper cannot. Copper heat sinks are machined and skived. Another method of manufacture is to solder the fins into the heat sink base. [0011] Another heat sink material that can be used is diamond. With a value of 2000 W/mK it exceeds that of copper by a factor of five. In contrast to metals, where heat is conducted by delocalized electrons, lattice vibrations are responsible for diamond's very high thermal conductivity. For thermal management applications, the outstanding thermal conductivity and diffusivity of diamond is an essential. CVD diamond may be used as a sub-mount for high-power integrated circuits and laser diodes. [0012] Composite materials can be used. Examples are a copper-tungsten pseudoalloy, AlSiC (silicon carbide in aluminum matrix), Dymalloy (diamond in copper-silver alloy matrix), and E-Material (beryllium oxide in beryllium matrix). Such materials are often used as substrates for chips, as their thermal expansion coefficient can be matched to ceramics and semiconductors. [0013] Fin efficiency is one of the parameters which makes a higher thermal conductivity material important. A fin of a heat sink may be considered to be a flat plate with heat flowing in one end and being dissipated into the surrounding fluid as it travels to the other. As heat flows through the fin, the combination of the thermal resistance of the heat sink impeding the flow and the heat lost due to convection, the temperature of the fin and, therefore, the heat transfer to the fluid, will decrease from the base to the end of the fin. This factor is called the fin efficiency and is defined as the actual heat transferred by the fin, divided by the heat transfer were the fin to be isothermal (hypothetically the fin having infinite thermal conductivity). Equations 5 and 6 are applicable for straight fins. [0000] η f = tanh  ( mL c ) mL c ( 5 ) mL c = 2   h f kt f  L f ( 6 ) [0014] Where: h f is the convection coefficient of the fin Air: 10 to 100 W/(m 2 K) Water: 500 to 10,000 W/(m 2 K) k is the thermal conductivity of the fin material Aluminum: 120 to 240 W/(m·K) L f is the fin height (m) t f is the fin thickness (m) [0022] Another parameter that concerns the thermal conductivity of the heat sink material is spreading resistance. Spreading resistance occurs when thermal energy is transferred from a small area to a larger area in a substance with finite thermal conductivity. In a heat sink, this means that heat does not distribute uniformly through the heat sink base. The spreading resistance phenomenon is shown by how the heat travels from the heat source location and causes a large temperature gradient between the heat source and the edges of the heat sink. This means that some fins are at a lower temperature than if the heat source were uniform across the base of the heat sink. This non-uniformity increases the heat sink's effective thermal resistance. [0023] A pin fin heat sink is a heat sink that has pins that extend from its base. The pins can be, for example, cylindrical, elliptical or square. A second type of heat sink fin arrangement is the straight fin. These run the entire length of the heat sink. A variation on the straight fin heat sink is a cross cut heat sink. A straight fin heat sink is cut at regular intervals but at a coarser pitch than a pin fin type. [0024] In general, the more surface area a heat sink has, the better it works. However, this is not always true. The concept of a pin fin heat sink is to try to pack as much surface area into a given volume as possible. As well, it works well in any orientation. Kordyban has compared the performance of a pin fin and a straight fin heat sink of similar dimensions. Although the pin fin has 194 cm 2 surface area while the straight fin has 58 cm 2 , the temperature difference between the heat sink base and the ambient air for the pin fin is 50° C. For the straight fin it was 44° C. or 6° C. better than the pin fin. Pin fin heat sink performance is significantly better than straight fins when used in their intended application where the fluid flows axially along the pins rather than only tangentially across the pins. [0025] Another configuration is the flared fin heat sink; its fins are not parallel to each other, but rather diverge with increasing distance from the base. Flaring the fins decreases flow resistance and makes more air go through the heat sink fin channel; otherwise, more air would bypass the fins. Slanting them keeps the overall dimensions the same, but offers longer fins. Forghan, et al. have published data on tests conducted on pin fin, straight fin and flared fin heat sinks. They found that for low approach air velocity, typically around 1 m/s, the thermal performance is at least 20% better than straight fin heat sinks. Lasance and Eggink also found that for the bypass configurations that they tested, the flared heat sink performed better than the other heat sinks tested. [0026] The heat transfer from the heatsink is mediated by two effects: conduction via the coolant, and thermal radiation. The surface of the heatsink influences its emissivity; shiny metal absorbs and radiates only a small amount of heat, while matte black radiates highly. In coolant-mediated heat transfer, the contribution of radiation is generally small. A layer of coating on the heatsink can then be counterproductive, as its thermal resistance can impair heat flow from the fins to the coolant. Finned heatsinks with convective or forced flow will not benefit significantly from being colored. In situations with significant contribution of radiative cooling, e.g. in case of a flat non-finned panel acting as a heatsink with low airflow, the heatsink surface finish can play an important role. Matte-black surfaces will radiate much more efficiently than shiny bare metal. The importance of radiative vs. coolant-mediated heat transfer increases in situations with low ambient air pressure (e.g. high-altitude operations) or in vacuum (e.g. satellites in space). [0027] Fourier, J. B., 1822, Theorie analytique de la chaleur, Paris; Freeman, A., 1955, translation, Dover Publications, Inc, N.Y. [0028] Kordyban, T., 1998, Hot air rises and heat sinks—Everything you know about cooling electronics is wrong, ASME Press, N.Y. [0029] Anon, Unknown, “Heat sink selection”, Mechanical engineering department, San Jose State University [27 Jan. 2010]. [0000] www.engr.sjsu.edu/ndej ong/ME%20146% 20files/Heat% 20Sink.pptwww.engr.sjsu.edu/ndejong/ME%20146%20files/Heat%20Sink.ppt [0030] Sergent, J. and Krum, A., 1998, Thermal management handbook for electronic assemblies, First Edition, McGraw-Hill. [0031] Incropera, F. P. and DeWitt, D. P., 1985, Introduction to heat transfer, John Wiley and sons, N.Y. [0032] Forghan, F., Goldthwaite, D., Ulinski, M., Metghalchi, M., 2001, Experimental and Theoretical Investigation of Thermal Performance of Heat Sinks, ISME May. [0033] Lasance, C. J. M and Eggink, H. J., 2001, A Method to Rank Heat Sinks in Practice: The Heat Sink Performance Tester, 21st IEEE SEMI-THERM Symposium. [0034] ludens.cl/Electron/Thermal.html [0035] Lienard, J. H., IV & V, 2004, A Heat Transfer Textbook, Third edition, MIT [0036] Saint-Gobain, 2004, “Thermal management solutions for electronic equipment” 22 Jul. 2008 www.fff.saint-gobain.com/Media/Documents/50000000000000001036/ThermaCool%20Brochure.pdf [0037] Jeggels, Y. U., Dobson, R. T., Jeggels, D. H., Comparison of the cooling performance between heat pipe and aluminium conductors for electronic equipment enclosures, Proceedings of the 14th International Heat Pipe Conference, Florianópolis, Brazil, 2007. [0038] Prstic, S., Iyengar, M., and Bar-Cohen, A., 2000, Bypass effect in high performance heat sinks, Proceedings of the International Thermal Science Seminar Bled, Slovenia, June 11-14. [0039] Mills, A. F., 1999, Heat transfer, Second edition, Prentice Hall. [0040] Potter, C. M. and Wiggert, D. C., 2002, Mechanics of fluid, Third Edition, Brooks/Cole. [0041] White, F. M., 1999, Fluid mechanics, Fourth edition, McGraw-Hill International. [0042] Azar, A, et al., 2009, “Heat sink testing methods and common oversights”, Qpedia Thermal E-Magazine, January 2009 Issue. [0000] www.qats.com/cpanel/UploadedPdf/January20092.pdf [0043] Several structurally complex heatsink designs are discussed in Hernon, US App. 2009/0321045, incorporated herein by reference. [0044] Heatsinks operate by removing heat from an object to be cooled into the surrounding air, gas or liquid through convection and radiation. Convection occurs when heat is either carried passively from one point to another by fluid motion (forced convection) or when heat itself causes fluid motion (free convection). When forced convection and free convection occur together, the process is termed mixed convection. Radiation occurs when energy, for example in the form of heat, travels through a medium or through space and is ultimately absorbed by another body. Thermal radiation is the process by which the surface of an object radiates its thermal energy in the form of electromagnetic waves. Infrared radiation from a common household radiator or electric heater is an example of thermal radiation, as is the heat and light (IR and visible EM waves) emitted by a glowing incandescent light bulb. Thermal radiation is generated when heat from the movement of charged particles within atoms is converted to electromagnetic radiation. [0045] A heatsink tends to decrease the maximum temperature of the exposed surface, because the power is transferred to a larger volume. This leads to a possibility of diminishing return on larger heatsinks, since the radiative and convective dissipation tends to be related to the temperature differential between the heatsink surface and the external medium. Therefore, if the heatsink is oversized, the efficiency of heat shedding is poor. If the heatsink is undersized, the object may be insufficiently cooled, the surface of the heatsink dangerously hot, and the heat shedding not much greater than the object itself absent the heatsink. [0046] The relationship between friction and convention in heatsinks is discussed by Frigus Primore in “A Method for Comparing Heat Sinks Based on Reynolds Analogy,” available at www.frigprim.com/downloads/Reynolds_analogy_heatsinks.PDF, last accessed Apr. 28, 2010. This article notes that for, plates, parallel plates, and cylinders to be cooled, it is necessary for the velocity of the surrounding fluid to be low in order to minimize mechanical power losses. However, larger surface flow velocities will increase the heat transfer efficiency, especially where the flow near the surface is turbulent, and substantially disrupts a stagnant surface boundary layer. Primore also discusses heatsink fin shapes and notes that no fin shape offers any heat dissipation or weight advantage compared with planar fins, and that straight fins minimize pressure losses while maximizing heat flow. Therefore, the art generally teaches that generally flat and planar surfaces are appropriate for most heatsinks. [0047] Frigus Primore, “Natural Convection and Inclined Parallel Plates,” www.frigprim.com/articels2/parallel_pI_Inc.html, last accessed Apr. 29, 2010, discusses the use of natural convection (i.e., convection due to the thermal expansion of a gas surrounding a solid heatsink in normal operating conditions) to cool electronics. One of the design goals of various heatsinks is to increase the rate of natural convection. Primore suggests using parallel plates to attain this result. Once again, Primore notes that parallel plate heat sinks are the most efficient and attempts to define the optimal spacing and angle (relative to the direction of the fluid flow) of the heat sinks according to the equations in FIG. 1 . [0048] In another article titled “Natural Convection and Chimneys,” available at www.frigprim.com/articels2/parallel_plchim.html, last accessed Apr. 29, 2010, Frigus Primore discusses the use of parallel plates in chimney heat sinks. One purpose of this type of design is to combine more efficient natural convection with a chimney. Primore notes that the design suffers if there is laminar flow (which creates a re-circulation region in the fluid outlet, thereby completely eliminating the benefit of the chimney) but benefits if there is turbulent flow which allows heat to travel from the parallel plates into the chimney and surrounding fluid. [0049] In “Sub-Grid Turbulence Modeling for Unsteady Flow with Acoustic Resonance,” available at www.metacomptech.com/cfd++/00-0473.pdf, last accessed Apr. 29, 2010, incorporated herein by reference, Paul Batten et al discuss that when a fluid is flowing around an obstacle, localized geometric features, such as concave regions or cavities, create pockets of separated flow which can generate self-sustaining oscillations and acoustic resonance. The concave regions or cavities serve to substantially reduce narrow band acoustic resonance as compared to flat surfaces. This is beneficial to a heat sink in a turbulent flow environment because it allows for the reduction of oscillations and acoustic resonance, and therefore for an increase in the energy available for heat transfer. [0050] In S. Liu, “Heat Transfer and Pressure Drop in Fractal Microchannel Heat Sink for Cooling of Electronic Chips,” 44 Heat Mass Transfer 221 (2007), Liu et al discuss a heatsink with a “fractal-like branching flow network.” Liu's heatsink includes channels through which fluids would flow in order to exchange heat with the heatsink. [0051] Y. J. Lee, “Enhanced Microchannel Heat Sinks Using Oblique Fins,” IPACK 2009-89059, similarly discusses a heat sink comprising a “fractal-shaped microchannel based on the fractal pattern of mammalian circulatory and respiratory system.” Lee's idea, similar to that of Liu, is that there would be channels inside the heatsink through which a fluid could flow to exchange heat with the heatsink. The stated improvement in Lee's heatsink is (1) the disruption of the thermal boundary layer development; and (2) the generation of secondary flows. [0052] Pence, D. V., 2002, “Reduced Pumping Power and Wall Temperature in Microchannel Heat Sinks with Fractal-like Branching Channel Networks”, Microscale Thermophys. Eng. 5, pp. 293-311, similarly, mentions heatsinks that have fractal-like channels allowing fluid to enter into the heat sink. The described advantage of Pence's structure is increased exposure of the heat sink to the fluid and lower pressure drops of the fluid while in the heatsink. [0053] In general, a properly designed heatsink system will take advantage of thermally induced convection or forced air (e.g., a fan). In general, a turbulent flow near the surface of the heatsink disturbs a stagnant surface layer, and improves performance. In many cases, the heatsink operates in a non-ideal environment subject to dust or oil; therefore, the heatsink design must accommodate the typical operating conditions, in addition to the as-manufactured state. [0054] Prior art heatsink designs have traditionally concentrated on geometry that is Euclidian, involving structures such as the pin fins, straight fins, and flares discussed above. [0055] N J Ryan, D A Stone, “Application of the FD-TD method to modelling the electromagnetic radiation from heatsinks”, IEEE International Conference on Electromagnetic Compatibility, 1997. 10th (1-3 Sep. 1997), pp: 119-124, discloses a fractal antenna which also serves as a heatsink in a radio frequency transmitter. [0056] Lance Covert, Jenshan Lin, Dan Janning, Thomas Dalrymple, “5.8 GHz orientation-specific extruded-fin heatsink antennas for 3D RF system integration”, 23 Apr. 2008 DOI: 10.1002/mop.23478, Microwave and Optical Technology Letters Volume 50, Issue 7, pages 1826-1831, July 2008 also provide a heatsink which can be used as an antenna. SUMMARY OF THE INVENTION [0057] Most heatsinks are designed using a linear or exponential relationship of the heat transfer and dissipating elements. A known geometry which has not generally been employed is fractal geometry. Some fractals are random fractals, which are also termed chaotic or Brownian fractals and include random noise components. In deterministic fractal geometry, a self-similar structure results from the repetition of a design or motif (or “generator”) using a recursive algorithm, on a series of different size scales. As a result, certain types of fractal images or structures appear to have self-similarity over a broad range of scales. On the other hand, no two ranges within the design are identical. [0058] A fractal is defined as “a rough or fragmented geometric shape that can be split into parts, each of which is (at least approximately) a reduced-size copy of the whole.” Mandelbrot, B. B. (1982). That is, there is a recursive algorithm which describes the structure. The Fractal Geometry of Nature. W. H. Freeman and Company. ISBN 0-7167-1186-9. This property is termed “self-similarity.” For a more detailed discussion of fractals, see the Wikipedia article thereon at en.wikipedia.org/wiki/Fractal (last accessed Apr. 14, 2010) incorporated herein by reference. Exemplary images of well-known fractal designs can also be viewed on the Wikipedia page. Due to the fact that fractals involve largely self-repeating patterns, each of which serves to increase the surface area in three-dimensional fractals (perimeter in two-dimensional fractals), three dimensional fractals in theory are characterized by infinite surface area (and two-dimensional fractals are characterized by infinite perimeter). In practical implementations, the scale of the smallest features which remain true to the generating algorithm may be 3-25 iterations of the algorithm. Less than three recursions, and the fractal nature is not apparent, while present manufacturing technologies limit the manufacture of objects with a large range of feature scales. [0059] This fractal nature is useful in a heatsink because the rate at which heat is transferred from a surface, either through convection or through radiation, is typically related to, and increasing with, the surface area. Of course, due to limitations in the technology used to build these heatsinks, engineering compromise is expected. However a feature of an embodiment of the designs proposed herein is that vortices induced by fluid flow over a heat transfer surface will be chaotically distributed over various elements of the surface, thus disrupting the stagnant surface boundary layer and increasing the effective surface area available for heat transfer, while avoiding acoustic resonance which may be apparent from a regular array of structures which produce vortices and turbulence. [0060] Further, a large physical surface area to volume ratio, which is generally useful in heatsink design, can still be obtained using the fractal model. In addition, fractal structures provide a plurality of concave regions or cavities, providing pockets of separated flow which can generate self-sustaining oscillations and acoustic resonance. These pockets serve to reduce the acoustic resonance in turbulent flowing fluid (as compared to flat or Euclidian surfaces), thus allowing for more effective heat transfer between the fractal structure and the surrounding fluid, thereby making the fractal structure ideal for a heatsink. [0061] U.S. Pat. No. 7,256,751, issued to Cohen, incorporated herein by reference, discusses fractal antennas. In the background of this patent, Cohen discusses Kraus' research, noting that Euclidian antennas with low area (and therefore low perimeter) exhibit very low radiation resistance and are thus inefficient. Cohen notes that the advantages of fractal antennas, over traditional antennas with Euclidian geometries, is that they can maintain the small area, while having a larger perimeter, allowing for a higher radiation resistance. Also, Cohen's fractal antenna features non-harmonic resonance frequencies, good bandwidth, high efficiency, and an acceptable standing wave ratio. [0062] In the instant invention, this same wave theory may be applied to fractal heatsinks, especially with respect to the interaction of the heat transfer fluid with the heatsink. Thus, while the heat conduction within a solid heatsink is typically not modeled as a wave (though modern thought applies phonon phenomena to graphene heat transport), the fluid surrounding the heating certainly is subject to wave phenomena, complex impedances, and indeed the chaotic nature of fluid eddies may interact with the chaotic surface configuration of the heatsink. [0063] The efficiency of capturing electric waves in a fractal antenna, achieved by Cohen, in some cases can be translated into an efficiency transferring heat out of an object to be cooled in a fractal heatsink as described herein. See, Boris Yakobson, “Acoustic waves may cool microelectronics”, Nano Letters, ACS (2010). Some physics scholars have suggested that heat can be modeled as a set of phonons. Convection and thermal radiation can therefore be modeled as the movement of phonons. A phonon is a quasiparticle characterized by the quantization of the modes of lattice vibration of solid crystal structures. Any vibration by a single phonon is in the normal mode of classical mechanics, meaning that the lattice oscillates in the same frequency. Any other arbitrary lattice vibration can be considered a superposition of these elementary vibrations. Under the phonon model, heat travels in waves, with a wavelength on the order of 1 μm. In most materials, the phonons are incoherent, and therefore a macroscopic scales, the wave nature of heat transport is not apparent or exploitable. [0064] The thermodynamic properties of a solid are directly related to its phonon structure. The entire set of all possible phonons combine in what is known as the phonon density of states which determines the heat capacity of a crystal. At absolute zero temperature (0 Kelvin or −273 Celsius), a crystal lattice lies in its ground state, and contains no phonons. A lattice at a non-zero temperature has an energy that is not constant, but fluctuates randomly about some mean value. These energy fluctuations are caused by random lattice vibrations, which can be viewed as a gas-like structure of phonons or thermal phonons. However, unlike the atoms which make up an ordinary gas, thermal phonons can be created and destroyed by random energy fluctuations. In the language of statistical mechanics this means that the chemical potential for adding a phonon is zero. For a more detailed description of phonon theory, see the Wikipedia article thereon available at en.wikipedia.org/wiki/Phonon (last accessed Apr. 16, 2010) incorporated herein by reference. [0065] In certain materials, such as graphene, phonon transport phenomena are apparent at macroscopic levels, which make phonon impedance measurable and useful. Thus, if a graphene sheet were formed to resonate at a particular phonon wavelength, the resonant energy would not be emitted. On the other hand, if the graphene sheet were configured using a fractal geometry, the phonon impedance would be well controlled over a broad range of wavelengths, with sharp resonances at none, leading to an efficient energy dissipation device. [0066] Many fractal designs are characterized by concave regions or cavities. See, for example, FIGS. 2 and 3 . While sets of concavities may be useful in improving aerodynamics and fluid dynamics to increase turbulence, if they are disposed in a regular array, they will likely produce an acoustic resonance, and may have peaks in a fluid impedance function. On the other hand, the multiscale nature of a fractal geometric design will allow the system to benefit from the concavities, while avoiding a narrowly tuned system. [0067] The present system proposes a fractal-shaped heatsink for the purpose of dissipating heat. The benefits of a fractal heatsink, over a traditional heatsink having a Euclidian geometry may include: (1) the fractal heatsink has a greater surface area, allowing for more exposure of the hot device to the surrounding air or liquid and faster dissipation of heat; and (2) due to the plethora of concave structures or cavities in fractal structures, the fractal heatsink is better able to take advantage of flow mechanics than a traditional heatsink, resulting in heat entering and exiting the heatsink more quickly (3) acoustic properties, especially in forced convection systems. [0068] The invention provides a heatsink to cool an object through convection or radiation. For the smallest heatsink elements, on the order of 10-100 nm, the focus of the heat transfer will be on radiation rather than convection. Electron emission and ionization may also be relevant. Larger heatsink elements, approximately >1 mm in size, will generally rely on convection as the primary form of heat transfer. [0069] In one embodiment, the heatsink comprises a heat exchange device with a plurality of heat exchange elements having a fractal variation therebetween. A heat transfer fluid, such as air, water, or another gas or liquid, is induced to flow through the heat exchange device. The heat transfer fluid has turbulent portions. The fractal variation in the plurality of heat exchange elements substantially reduces the narrow band acoustic resonance as compared to a heatsink having a linear or Euclidian geometric variation between the plurality heat exchange elements. The turbulent flow also disturbs the stagnant surface boundary layer, leading to more efficient heat transfer. [0070] When a heat transfer fluid (air, gas or liquid) is induced to flow over a surface, there may be turbulence in the fluid. The fractal shape of the heatsink serves to reduce the energy lost due to the turbulence, while increasing the surface area of the heatsink subject to turbulence, due to the plethora of concave regions, cavities, and pockets. Therefore, the efficiency of heat transfer may be increased as compared to a heat exchange device having a linear or Euclidian geometric variation between several heat exchange elements. [0071] Preferably, the heat exchange device will include a highly conductive substance whose heat conductivity exceeds 850 W/(m*K). Examples of such superconductors include graphene, diamond, and diamond-like coatings. Alternatively, the heat exchange device may include carbon nanotubes. [0072] Various variations on this heatsink will be apparent to skilled persons in the art. For example, the heatsink could include a heat transfer surface that is connected to the heat exchange device and is designed to accept a solid to be cooled. Alternatively, there could be a connector that is designed to connect with a solid to be cooled in at least one point. In another embodiment, there are at least three connectors serving to keep the solid and the heatsink in a fixed position relative to one another. Various connectors will be apparent to persons skilled in the art. For example, the connector could be a point connector, a bus, a wire, a planar connector or a three-dimensional connector. In another embodiment, the heatsink has an aperture or void in the center thereof designed to accept a solid to be cooled. [0073] This heatsink is intended to be used to cool objects, and may be part of a passive or active system. Modern three-dimensional laser and liquid printers can create objects such as the heatsinks described herein with a resolution of features on the order of about 16 μm, making it feasible for those of skilled in the art to use such fabrication technologies to produce objects with a size below 10 cm. Alternatively, larger heatsinks, such as car radiators, can be manufactured in a traditional manner, designed with an architecture of elements having a fractal configuration. For example, a liquid-to-gas heat exchanger (radiator) may be provided in which segments of fluid flow conduit have a fractal relationship over three levels of recursion, i.e., paths with an average of at least two branches. Other fractal design concepts may be applied concurrently, as may be appropriate. [0074] Yet another embodiment of the invention involves a method of cooling a solid by connecting the solid with a heatsink. The heatsink comprises a heat exchange device having a plurality of heat exchange elements having a fractal variation therebetween. A heat transfer fluid having turbulent portions is induced to flow with respect to the plurality of heat exchange elements. The fractal variation in the plurality of heat exchange elements serves to substantially reduce narrow band resonance as compared to a corresponding heat exchange device having a linear or Euclidean geometric variation between a plurality of heat exchange elements. [0075] A preferred embodiment provides a surface of a solid heatsink, e.g., an internal or external surface, having fluid thermodynamical properties adapted to generate an asymmetric pattern of vortices over the surface over a range of fluid flow rates. For example, the range may comprise a range of natural convective fluid flow rates arising from use of the heatsink to cool a heat-emissive object. The range may also comprise a range of flow rates arising from a forced convective flow (e.g., a fan) over the heatsink. [0076] The heatsink may cool an unconstrained or uncontained fluid, generally over an external surface of a heatsink, or a constrained or contained fluid, generally within an internal surface of a heatsink. BRIEF DESCRIPTION OF THE DRAWINGS [0077] FIG. 1 shows a set of governing equations for a parallel plate heatsink. [0078] FIG. 2 illustrates a fractal heatsink that is an exemplary embodiment of the invention. In this embodiment, the heatsink is placed adjacent to the object to be cooled. [0079] FIG. 3 illustrates a fractal heatsink that is an exemplary embodiment of the invention. In this embodiment, the heatsink is placed either adjacent to or surrounding the object to be cooled. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS [0080] FIG. 2 illustrates a heatsink implementing an exemplary embodiment of this invention. Note that the illustration is in two dimensions, but a three dimensional embodiment is both possible and preferred. There is a heat transfer surface 100 that allows the heatsink to rest comfortably on a surface, such as the solid to be cooled 190 . In the illustrated embodiment, the heat transfer surface 100 is roughly planar, having a closed Euclidian cross-section on the bottom. However, it might also have another shape, for example if the solid to be cooled does not have a planar face. A fractal-shaped heat exchange device begins at point 110 . While only one fractal heatsink is illustrated here, skilled persons in the art will recognize other similar fractal heatsinks that are also intended to be covered by the invention. Note that the heatsink has three branches leaving from point 110 —branch 120 , branch 140 , and branch 160 . Also note that the branch structure initiating from point 110 is nearly identical to that at point 122 and 142 , even though only point 110 is a true starting point. Thus, the fractal property of self-similarity is preserved. We call the structure that begins at point 110 the “first motif,” the structure from point 122 the “second motif,” and the structure that begins from point 142 the “third motif.” Note that, in the embodiment illustrated in FIG. 2 , the replication from first to second motif and from second to third motif involves a linear displacement (upward) and a change of scale. In branches not going in the same direction as the prior branch, there is also a rotation. Under the limitations for ideal fractals, the second motif and third motif must be a smaller, exact copy of the first motif. However, due to the limitations imposed by human-made structures and machines, the fractals designed here are generally finite and the second motif will thus be an inexact copy of the first motif, i.e. if there are N levels starting from the first motif, the second motif level will have N-1 levels, if N is very large, the difference is insignificant. In other words, the self-similarity element required in fractals is not preserved perfectly in the preferred designs due to the limitations of available machinery. In addition, the benefits are achieved without requiring fractal relationships over more than a few “orders” of magnitude (iterations of the fractal recursive algorithm). For example, in the embodiment illustrated in FIG. 2 , there are no continuing branch divisions and iterations at point 162 , even though an ideal fractal would have them. In an ideal fractal, there would be an infinite number of sub-branches from 110 , 122 , and 142 . However, an imperfect fractal shape, as illustrated in FIG. 2 , will serve the purposes of this invention. [0081] Persons of ordinary skill in the art will appreciate the advantages offered by the structure 110 in FIG. 2 . The fractal heatsink has a much larger surface area than the heat transfer surface alone because all of the “branches” and “leaves” of the fern-like fractal shape serve to increase the surface area. In addition, if a heat transfer fluid is induced to flow above the heat transfer surface 100 , the turbulent portions of the heat transfer fluid near the surface will be increased by the textures inherent in the fractal variation in the heat exchange element 110 . Because the fractal patterns is itself non-identically repeating within the fractal design, this will serve to substantially reduce narrow band acoustic resonance as compared to a corresponding heat exchange device having a repeating design, e.g., a linear or geometric variation between several heat exchange elements, thereby further aiding in the heat transfer process. [0082] In a preferred embodiment, the heat transfer surface 100 and the roughly fractal-shaped heat exchange element 110 are all made out of an efficient heat conductor, such as copper or aluminum, or more preferably, having a portion whose heat conductivity exceeds 850 W/(m*K), such as graphene with a heat conductivity of between 4840 and 5300 W/(m*K) or diamond with a heat conductivity between 900 and 2320 W/(m*K). This would allow heat to quickly enter the heatsink from the solid and for heat to quickly exit the heatsink through the branches and leaves of the fern-like fractal 110 . In another embodiment, the heatsink is formed, at least in part, of carbon nanotubes, which display anisotropic heat conduction, with an efficient heat transfer along the long axis of the tube. Carbon nanotubes are submicroscopic hollow tubes made of a chicken-wire-like or lattice of carbon atoms. These tubes have a diameter of just a few nanometers and are highly heat conductive, transferring heat much faster than diamond, and in some cases comparable to graphene. See web.mit.edu/press/2010/thermopower-waves.html (last accessed Apr. 15, 2010) incorporated herein by reference. [0083] Also note that this exemplary embodiment provides a plethora of openings, e.g. 124 and 126 , between the branches or fractal subelements to ensure that all of the branches are exposed to the surrounding air, gas or liquid and to allow the heat to escape from the heatsink into the surroundings. In one embodiment of the invention, at least two of these openings are congruent, as are openings 124 and 126 illustrated here. An embodiment of the invention allows the openings to be filled with the air or liquid from the surrounding medium. Due to the limitation imposed by the solid's flat shape, it is not possible to increase the exposure of the fern-like fractal to the solid. However, the air or liquid outside of the solid are perfect for the fractal's exposure. [0084] Under the phonon model of heat exchange, applicable to carbon nanotubes, graphene materials, and perhaps others, the fractal shape is advantageous to ensure the escape of the phonons into the surrounding fluid medium because the fractal guarantees close to maximal surface exposure to the medium and does not have many parts that are not exposed, as is a problem with many prior art heatsinks. Skilled persons in the art will realize that this could be achieved through many known structures. For example, graphene, which is one-atom-thick carbon and highly heat conductive, would be an advantageous material to use to build the fractal heatsink herein described. [0085] When a turbulently flowing fluid passes around an obstacle, concave regions or cavities in the obstacle create pockets of separated flow which generates self-sustaining oscillations and acoustic resonance. The concave regions or cavities have substantially reduced narrow band acoustic resonance as compared to flat regions on the obstacle. This allows for more energy to be available for heat transfer. Skilled persons in the art will note that fractal structure 110 , as many other fractal structures, has a plurality of concave regions to allow for an implementation of this effect. [0086] FIG. 3 illustrates another embodiment of the invention. A solid to be cooled that has an arbitrary shape 290 is located inside (illustrated) or outside (not illustrated) a two-dimensional or three-dimensional roughly fractal shaped 210 heatsink. In one embodiment, the heatsink 210 has an aperture 270 designed to hold the solid. Alternatively, the solid to be cooled might be located outside of the heatsink (not illustrated). Note that, as in FIG. 2 , the fractal heat exchange element has multiple motifs, starting with the large triangle at 210 , to progressively smaller triangles at 220 and 230 . However, note that the fractal does not keep extending infinitely and there are no triangles smaller than the one at 230 . In other words, the fractal heatsink 210 has multiple recursive fractal iterations 220 and 230 , but the fractal iterations stop at level 230 for simplicity of design and manufacturability. Also note that the fractal submotifs 220 and 230 are of different dimensional sizes from the original fractal motif 210 and protrude from the original fractal shape 210 . Here, the first motif is a large triangle, and the latter motifs are smaller triangles, which involve a rotation, linear displacement, and change of scale of the prior motif. In one embodiment, the fractal shape has some apertures in it (not illustrated) to allow the solid to be cooled to connect with other elements. Also, the solid to be cooled is connected to the fractal shape at point connector 240 and through bus wires at 250 and 260 . The solid should be connected to the fractal heatsink in at least one point, either through a point connection, a bus wire connection, or some other connection. If it is desired that the solid be fixed inside the heatsink, there may be at least three connection points, as illustrated. However, only one connection point is necessary for heat convection and radiation from the solid to the heatsink. Preferably, the point or bus wire connection is built using a strong heat conductor, such as carbon nanotubes or a diamond-like coating. [0087] Note that, as in FIG. 1 , the fractal structure 210 in FIG. 2 has multiple concave regions or cavities. When a turbulently flowing fluid passes around this fractal heatsink, the concave regions or cavities substantially reduce the narrow band acoustic resonance as compared to a flat or Euclidian structure. This allows for more energy to be available to for heat transfer. [0088] In yet another embodiment of the invention, the heatsink 210 in FIG. 3 could be constructed without the connections at points 240 , 250 , and 260 . In one embodiment, a liquid or gas would fill the aperture 270 with the intent that the liquid or gas surround the solid to be cooled, hold it in place, or suspend it. Preferably, the liquid or gas surrounding the solid would conduct heat from the solid to the heatsink, which would then cause the heat to exit. [0089] Those skilled in the art will recognize many ways to fabricate the heatsinks described herein. For example, modern three-dimensional laser and liquid printers can create objects such as the heatsinks described herein with a resolution of features on the order of 16 μm. Also, it is possible to grow a crystal structure using a recursive growth algorithm or through crystal growth techniques. For example, US Patent Application No. 2006/0037177 by Blum, incorporated herein by reference, describes a method of controlling crystal growth to produce fractals or other structures through the use of spectral energy patterns by adjusting the temperature, pressure, and electromagnetic energy to which the crystal is exposed. This method might be used to fabricate the heatsinks described herein. For larger heatsinks, such as those intended to be used in car radiators, traditional manufacturing methods for large equipment can be adapted to create the fractal structures described herein. [0090] In this disclosure, we have described several embodiments of this broad invention. Persons skilled in the art will definitely have other ideas as to how the teachings of this specification can be used. It is not our intent to limit this broad invention to the embodiments described in the specification. Rather, the invention is limited by the following claims.

A heatsink comprising a heat exchange device having a plurality of heat exchange elements each having a surface boundary with respect to a heat transfer fluid, having a fractal variation therebetween, wherein the heat transfer fluid is induced to flow with respect to the plurality of fractally varying heat exchange elements such that flow-induced vortices are generated at non-corresponding locations of the plurality of fractally varying heat exchange elements, resulting in a reduced resonance as compared to a corresponding heat exchange device having a plurality of heat exchange elements that produce flow-induced vortices at corresponding locations on the plurality of heat exchange elements.

FIELD OF THE INVENTION [0001] The present invention relates to herbicidally active compositions, which comprise 5-ethyl-2-[(RS)-4-isopropyl-4-methyl-5-oxo-2-imidazolin-2-yl]nicotinic acid (common name: imazethapyr) and (5RS)-2-[(EZ)-1-(Ethoxyimino)butyl]-3-hydroxy-5-[(3RS)-thian-3-yl]cyclohex-2-en-1-one (common name: cycloxydim) and optionally 2-[(RS)-4-isopropyl-4-methyl-5-oxo-2-imidazolin-2-yl]-5-methoxymethylnicotinic acid (common name: imazamox). [0002] Further, the invention relates to herbicidally active compositions comprising imazethapyr, cycloxydim and bentazone and optionally imazamox. BACKGROUND OF THE INVENTION [0003] In crop protection, it is desirable in principle to increase the specificity and the reliability of the action of active compounds. In particular, it is desirable for the crop protection product to control the harmful plants effectively and, at the same time, to be tolerated by the useful plants in question. [0004] 5-ethyl-2-[(RS)-4-isopropyl-4-methyl-5-oxo-2-imidazolin-2-yl]nicotinic acid (common name: imazethapyr; formula I) is an active compound from the group of imidazolinone herbicides, which are known e.g. from Shaner, D. L., O'Conner, S. L, The Imidazolinone Herbicides, CRC Press Inc., Boca Raton, Fla. 1991 and also from The Compendium of Pesticide Common Names http://www.alanwood.net/pesticides/. [0000] [0005] (5RS)-2-[(EZ)-1-(Ethoxyimino)butyl]-3-hydroxy-5-[(3RS)-thian-3-yl]cyclohex-2-en-1-one (common name: cycloxydim; formula II) is an active compound from the group of cyclohexanedione herbicides. [0000] [0006] 2-[(RS)-4-isopropyl-4-methyl-5-oxo-2-imidazolin-2-yl]-5-methoxymethylnicotinic acid (common name: imazamox; formula III) is an active compound from the group of imidazolinone herbicides, which are known e.g. from Shaner, D. L. O'Conner, S. L The Imidazolinone Herbicides, CRC Press Inc., Boca Raton, Fla. 1991 and also from The Compendium of Pesticide Common Names http://www.alanwood.net/pesticides/. [0000] [0007] 3-isopropyl-1H-2,1,3-benzothiadiazin-4(3H)-one 2,2-dioxide (common name: bentazone; formula IV) is an active compound from the group of thiadiazine herbicides. [0000] [0008] Although both imazethapyr and combinations of imazethapyr with imazamox (for example commercialized under the Odyssey® brand) are highly effective pre- and post-emergence herbicides, in some cases they do not provide a sufficient control of the relevant harmful plants and their activity at low application rates is not always satisfactory. Apart from that, its compatibility with certain dicotyledonous crop plants such as soybean, peanuts or other pulse or leguminous crops is not always satisfactory, i.e. in addition to the harmful plants, the crop plants are also damaged to an extent which is not acceptable. Though it is in principle possible to spare crop plants by lowering the application rates, the extent of the control of harmful plants is naturally also reduced. DETAILED DESCRIPTION OF THE INVENTION [0009] It is an object of the present invention to provide herbicidal compositions, which show enhanced herbicide action against undesirable harmful plants, in particular against Acalypha species such as Acalypha indica, Dinebra species such as Dinebra Arabica, Cynotis spec such as Cynotis axillaris, Parthenium spec such as Parthenium hysterophorus, Physalis spec such as Physalis minima, Digera spec such as Digera arvensis, Alopecurus myosuroides, Apera spicaventi, Brachiaria spec. such as Brachiaria deflexa or Brachiaria plantaginea, Echinochloa spec. such as Echinochloa colonum, Leptochloa spec. such as Leptochloa fusca, Rottboellia cochinchinensis, Digitaria sanguinalis, Eleusine indica, Saccharum spontaneum, Cynodon dactylon, Euphorbia hirta, Euphorbia geniculata, Commelina benghalensis, Commelina communis, certain undesired Oryza spec. such as weedy rice or red rice ( Oryza sativa ), Phalaris spec. such as Phalaris canariensis, Celosia argentea, Xanthium strumarium, Papaver rhoeas, Geranium spec, Brassica spec, Avena fatua, Bromus spec., Lolium spec., Phalan spec., Setaria spec., Digitaria spec., brachiaria spec., Amaranthus spec., Chenopodium spec., Abutilon theophrasti, Galium aparine, Veronica spec., or Solanum spec. and/or to improve their compatibility with crop plants, such as soybean, peanut, pea, bean, lentil, green gram, black gram, cluster bean, fenugreek, other pulse or leguminous crops, or crops which are tolerant to the action of acetohydroxyacid synthase inhibiting herbicides, such as for example Clearfield® wheat, Clearfield® barley, Clearfield® corn, Clearfield® lentil, Clearfield® oilseed rape or canola, Clearfield® rice, Cultivance® soybean and/or Clearfield® sunflower. The composition should also have a good pre-emergence herbicidal activity. [0010] We consider this object to be achievable, by herbicidally active compositions comprising imazethapyr or agriculturally acceptable salts thereof and cycloxydim or agriculturally acceptable salts thereof and optionally imazamox or agriculturally acceptable salts thereof. [0011] Further, we consider this object to be achievable, by herbicidally active compositions comprising imazethapyr or agriculturally acceptable salts thereof, cycloxydim or agriculturally acceptable salts thereof and bentazone or agriculturally acceptable salts thereof and optionally imazamox or agriculturally acceptable salts thereof. [0012] The invention relates in particular to compositions in the form of agriculturally acceptable compositions as defined above. [0013] The invention furthermore relates to the use of compositions as defined herein for controlling undesirable vegetation in crops. When using the compositions of the invention for this purpose the imazethapyr agriculturally acceptable salts thereof and cycloxydim agriculturally acceptable salts thereof and optionally imazamox or agriculturally acceptable salts thereof can be applied simultaneously or in succession in crops, where undesirable vegetation may occur. [0014] When using the compositions of the invention for this purpose the imazamox or agriculturally acceptable salts thereof and cycloxydim or agriculturally acceptable salts thereof can be applied simultaneously or in succession in crops, where undesirable vegetation may occur. [0015] The invention furthermore relates to the use of compositions as defined herein for controlling undesirable vegetation in crops which, by genetic engineering or by breeding, are tolerant to one or more herbicides, e.g. acetohydroxyacid synthase inhibiting herbicides such as imazethapyr or imazamox, and/or pathogens such as plant-pathogenous fungi, and/or to attack by insects; preferably tolerant to one or more herbicides that act as acetohydroxyacid synthase inhibitors. [0016] The invention furthermore relates to a method for controlling undesirable vegetation, which comprises applying an herbicidal composition according to the present invention to the undesirable plants. Application can be done before, during and/or after, preferably during and/or after, the emergence of the undesirable plants. [0017] The combination of imazethapyr or agriculturally acceptable salts thereof and cycloxydim or agriculturally acceptable salts thereof and optionally imazamox or agriculturally acceptable salts thereof can be applied simultaneously or in succession. The combination of imazethapyr or agriculturally acceptable salts thereof, cycloxydim or agriculturally acceptable salts thereof and bentazone or agriculturally acceptable salts thereof and optionally imazamox or agriculturally acceptable salts thereof can be applied simultaneously or in succession. [0018] The invention in particular relates to a method for controlling undesirable vegetation in crops, which comprises applying an herbicidal composition according to the present invention in crops where undesirable vegetation occurs or might occur. [0019] The invention furthermore relates to a method for controlling undesirable vegetation, which comprises allowing a composition according to the present invention to act on plants, their habitat or on seed. [0020] In the methods of the present invention it is immaterial whether the combination of imazethapyr or agriculturally acceptable salts thereof and cycloxydim or agriculturally acceptable salts thereof and optionally imazamox or agriculturally acceptable salts thereof or the combination imazethapyr or agriculturally acceptable salts thereof, cycloxydim or agriculturally acceptable salts thereof and bentazone or agriculturally acceptable salts thereof and optionally imazamox or agriculturally acceptable salts thereof are formulated and applied jointly or separately, and, in the case of separate application, in which order the application takes place. It is only necessary, that the combination of imazethapyr or agriculturally acceptable salts thereof and cycloxydim or agriculturally acceptable salts thereof and optionally imazamox or agriculturally acceptable salts thereof or the combination of imazethapyr or agriculturally acceptable salts thereof, cycloxydim or agriculturally acceptable salts thereof, bentazone or agriculturally acceptable salts thereof and optionally imazamox or agriculturally acceptable salts thereof are applied in a time frame, which allows simultaneous action of the active ingredients on the plants. [0021] The invention also relates to an herbicide formulation, which comprises an agriculturally acceptable composition as defined herein and at least one carrier material, including liquid and/or solid carrier materials. [0022] It is believed that the compositions according to the invention have better herbicidal activity against harmful plants than would have been expected by the herbicidal activity of the individual compounds. In other words, the joint action of the imazethapyr and cycloxydim combination, the imazethapyr, cycloxydim and imazamox combination, the imazethapyr, cycloxydim and bentazone combination and the imazethapyr, cycloxydim, bentazone and imazamox combination result in an enhanced activity against harmful plants in the sense of a synergy effect (synergism). For this reason, the compositions can, based on the individual components, be used at lower application rates to achieve a herbicidal effect comparable to the individual components. The compositions of the invention also show an accelerated action on harmful plants, i.e. damaging of the harmful plants is achieved more quickly in comparison with application of the individual herbicides. Moreover, the compositions of the present invention provide good pre-emergence herbicidal activity, i.e. the compositions are particularly useful for combating/controlling harmful plants before their emergence. Apart from that, the compositions of the present invention show good crop compatibility, i.e. their use in crops leads to a reduced damage of the crop plants. [0023] As used herein, the terms “controlling” and “combating” are synonyms. [0024] As used herein, the terms “undesirable vegetation” and “harmful plants” are synonyms. [0025] Where reference is made to imazethapyr, this includes the free acid as well as agriculturally acceptable salts thereof. [0026] Examples of agriculturally acceptable salts of imazethapyr include alkaline or earth alkaline metals or ammonium or organoammonium salts, for instance, sodium, potassium, ammonium and isopropyl ammonium. Particularly preferred is imazethapyr ammonium salt. Reference to a salt includes the anhydrous form as well as hydrated forms thereof. [0027] Where reference is made to cycloxydim, this includes the free acid as well as agriculturally acceptable salts thereof. [0028] Examples of agriculturally acceptable salts of cycloxydim include alkaline or earth alkaline metals or ammonium or organoammonium salts, for instance, lithium, sodium, potassium, ammonium and isopropyl ammonium. Particularly preferred are the lithium and the sodium salt of cycloxydim. Reference to a salt includes the anhydrous form as well as hydrated forms thereof. [0029] Where reference is made to imazamox, this includes the free acid as well as agriculturally acceptable salts thereof. [0030] Examples of agriculturally acceptable salts of imazamox include alkaline or earth alkaline metals or ammonium or organoammonium salts, for instance, sodium, potassium, ammonium and isopropyl ammonium. Particularly preferred is imazamox ammonium salt. Reference to a salt includes the anhydrous form as well as hydrated forms thereof. [0031] Where reference is made to bentazone, this includes the free base as well as agriculturally acceptable salts thereof. [0032] Examples of agriculturally acceptable salts of bentazone include alkaline or earth alkaline metals or ammonium or organoammonium salts, for instance, lithium, sodium, potassium, magnesium, calcium, ammonium, methylammonium, dimethylammonium, methyltrioctylammonium and isopropyl ammonium as well as the diethanolamine and choline salts. Particularly preferred is bentazone sodium salt. Reference to a salt includes the anhydrous form as well as hydrated forms thereof. [0033] The compositions of the invention comprise the herbicidally active agents imazethapyr and cycloxydim and optionally imazamox. [0034] Imazethapyr may be present in the form of its racemate or in the form of the pure R- or S-enantiomers (including salts as defined above). In particular suitable is R-imazethapyr. Cycloxydim may be present in the form of its racemate or in the form of the pure R- or S-enantiomers, or it may be present in the form of pure stereoisomers with respect to the two chiral centers and the ethoxyimino group, or as mixtures of two or more such stereoisomers (in each case, including salts as defined above). Particularly preferred is a racemic mixture of stereoisomers. Imazamox may be present in the form of its racemate or in the form of the pure R- or S-enantiomers (including salts as defined above). In particular suitable is R-imazamox. [0035] The compositions of the invention may also comprise one or more safeners. Safeners, also termed as herbicide safeners, are organic compounds which in some cases lead to better crop plant compatibility when applied jointly with specifically acting herbicides. Some safeners are themselves herbicidally active. In these cases, the safeners act as antidote or antagonist in the crop plants and thus reduce or even prevent damage to the crop plants. However, in the compositions of the present invention, safeners are generally not required. Therefore, a preferred embodiment of the invention relates to compositions which contain no safener or virtually no safener (i.e. less than 1% by weight, based on the total amount of imazethapyr and cycloxydim and optionally imazamox). [0036] Suitable safeners, which can be used in the compositions according to the present invention are known in the art, e.g. from The Compendium of Pesticide Common Names (http://www.alanwood.net/pesticides/); Farm Chemicals Handbook 2000 Vol. 86, Meister Publishing Company, 2000; B. Hock, C. Fedtke, R. R. Schmidt, Herbizide, Georg Thieme Verlag, Stuttgart 1995; W. H. Ahrens, Herbicide Handbook, 7th Edition, Weed Science Society of America, 1994; and K. K. Hatzios, Herbicide Handbook, Supplement to 7th Edition, Weed Science Society of America, 1998. [0037] Safeners include benoxacor, cloquintocet, cyometrinil, cyprosulfamide, dichlormid, dicyclonon, dietholate, fenchlorazole, fenclorim, flurazole, fluxofenim, furilazole, isoxadifen, mefenpyr, mephenate, naphthalic anhydride, 2,2,5-trimethyl-3-(dichloracetyl)-1,3-oxazolidine, 4-(dichloroacetyl)-1-oxa-4-azaspiro[4.5]decane and oxabetrinil, as well as thereof agriculturally acceptable salts and, provided they have a carboxyl group, their agriculturally acceptable derivatives. 2,2,5-Trimethyl-3-(dichloroacetyl)-1,3-oxazolidine [CAS No. 52836-31-4] is also known under the name R-29148.4-(Dichloroacetyl)-1-oxa-4-azaspiro[4.5]decane [CAS No. 71526-07-03] is also known under the names AD-67 and MON 4660. [0038] As safener, the compositions according to the invention particularly preferably com-prise at least one of the compounds selected from the group of benoxacor, cloquintocet, cyprosulfamide, dichlormid, fenchlorazole, fenclorim, fluxofenim, furilazole, isoxadifen, mefenpyr, naphthalic anhydride, 2,2,5-trimethyl-3-(dichloroacetyl)-1,3-oxazolidine, and 4-(dichloroacetyl)-1-oxa-4-azaspiro[4.5]decane and oxabetrinil; and the agriculturally acceptable salt thereof and, in the case of compounds having a COOH group, an agriculturally acceptable derivative as defined below. [0039] A preferred embodiment of the invention relates to compositions which contain no safener or virtually no safener (i.e. less than 1% by weight, based on the total amount of imazethapyr and cycloxydim and optionally imazamox) is applied. [0040] If the compounds of herbicidally active compounds mentioned as safeners (see above) have functional groups, which can be ionized, they can also be used in the form of their agriculturally acceptable salts. In general, the salts of those cations are suitable whose cations have no adverse effect on the action of the active compounds (“agricultural acceptable”). [0041] In general, the salts of those cations are suitable whose cations have no adverse effect on the action of the active compounds (“agriculturally acceptable”). Preferred cations are the ions of the alkali metals, preferably of lithium, sodium and potassium, of the alkaline earth metals, preferably of calcium and magnesium, and of the transition metals, preferably of manganese, copper, zinc and iron, furthermore ammonium and substituted ammonium (hereinafter also termed as organoammonium) in which one to four hydrogen atoms are replaced by C1-C4-alkyl, hydroxy-C1-C4-alkyl, C1-C4-alkoxy-C1-C4-alkyl, hydroxy-C1-C4-alkoxy-C1-C4-alkyl, phenyl or benzyl, preferably ammonium, methylammonium, isopropylammonium, dimethylammonium, diisopropylammonium, trimethylammonium, tetramethylammonium, tetraethylammonium, tetrabutylammonium, 2-hydroxyethylammonium, 2-(2-hydroxyethoxy)eth-1-ylammonium, di(2-hydroxyeth-1-yl)ammonium, benzyltrimethylammonium, benzyltriethylammonium, furthermore phosphonium ions, sulfonium ions, preferably tri(C 1 -C 4 -alkyl)sulfonium such as trimethylsulfonium, and sulfoxonium ions, preferably tri(C 1 -C 4 -alkyl)sulfoxonium. [0042] In the compositions according to the invention, the compounds that carry a carboxyl group can also be employed in the form of agriculturally acceptable derivatives, for example as amides such as mono- or di-C 1 -C 6 -alkylamides or arylamides, as esters, for example as allyl esters, propargyl esters, C 1 -C 10 -alkyl esters or alkoxyalkyl esters, and also as thioesters, for example as C 1 -C 10 -alkyl thioesters. Preferred mono- and di-C 1 -C 6 -alkylamides are the methyl- and the dimethylamides. Preferred arylamides are, for example, the anilidines and the 2-chloroanilides. Preferred alkyl esters are, for example, the methyl, ethyl, propyl, isopropyl, butyl, isobutyl, pentyl, mexyl (1-methylhexyl) or isooctyl (2-ethylhexyl) esters. Preferred C 1 -C 4 -alkoxy-C 1 -C 4 -alkyl esters are the straight-chain or branched C 1 -C 4 -alkoxyethyl esters, for example the methoxyethyl, ethoxyethyl or butoxyethyl esters. An example of the straight-chain or branched C 1 -C 10 -alkyl thioesters is the ethyl thioester. Preferred derivatives are the esters. [0043] The compositions of the present invention are suitable for controlling a large number of harmful plants, including monocotyledonous weeds, in particular annual weeds such as gramineous weeds (grasses) including Echinochloa species such as barnyardgrass ( Echinochloa crusgalli var crus - galli ) or Echinochloa colonum, Digitaria species such as crabgrass ( Digitaria sanguinalis ), certain undesired Oryza spec. such as weedy rice or red rice ( Oryza sativa ), Setaria species such as green foxtail ( Setaria viridis ) and giant foxtail ( Setaria faberil ), Sorghum species such as johnsongrass ( Sorghum halepense Pers. ), Avena species such as wild oats ( Avena fatua ), Cenchrus species such as Cenchrus echinatus, Bromus species, Lolium species, Leptochloa spec. such as Leptochloa fusca, Phalaris species such as Phalaris canariensis, Eriochloa species, Panicum species, Brachiaria species such as Brachiaria deflexa or Brachiaria plantaginea, annual bluegrass ( Poa annua ), blackgrass ( Alopecurus myosuroides ), Aegilops cylindrica, Agropyron repens, Apera spica - venti, Eleusine indica, Cynodon dactylon, Rottboellia cochinchinensis, Dinebra species such as Dinebra arabica, Saccharum spontaneum, and the like. [0044] The invention is particular suitable for controlling monocotyledonous weeds selected from the group consisting of Echinochloa spec. such as Echinochloa colonum, Leptochloa spec. such as Leptochloa fusca, certain undesired Oryza spec. such as weedy rice or red rice ( Oryza sativa ), Phalaris spec. such as Phalaris canariensis, Brachiaria spec. such as Briachiaria deflexa or Brachiaria plataginea. [0045] The compositions of the present invention are also suitable for controlling a large number of dicotyledonous weeds, in particular broad leaf weeds including Acalypha species such as Acalypha indica, Polygonum species such as wild buckwheat ( Polygonum convolvolus ), Amaranthus species such as pigweed ( Amaranthus retroflexus ) or Amaranthus viridis, Chenopodium species such as common lambsquarters ( Chenopodium album L.), Sida species such as prickly sida ( Sida spinosa L.), Ambrosia species such as common ragweed ( Ambrosia artemisiifolia ), Acanthospermum species, Anthemis species, Atriplex species, Cirsium species, Convolvulus species, Conyza species, Cassia species, Commelina species such as for example Commelina communis or Commelina benghalens, Datura species, Euphorbia species such as for example Euphorbia hirta or Euphorbia geniculata, Geranium species, Galinsoga species, morningglory ( Ipomoea species), Lamium species, Malva species, Matricaria species, Sysimbrium species, Solanum species, Xanthium species such as for example Xanthium strumiarium, Veronica species, Viola species, common chickweed ( Stellaria media ), velvetleaf ( Abutilon theophrast ), Hemp sesbania ( Sesbania exaltata Cory), Anoda cristata, Bidens pilosa, Brassica kaber, Capsella bursa - pastoris, Centaurea cyanus, Galeopsis tetrahit, Galium aparine, Helianthus annuus, Desmodium tortuosum, Kochia scoparia, Mercurialis annua, Myosotis arvensis, Papaver rhoeas, Raphanus raphanistrum, Salsola kali, Sinapis arvensis, Sonchus arvensis, Thlaspi arvense, Tagetes minuta, Richardia braslliensis, Rumex crispus, Rumex obtusifolius, Heracleaum sphondylium, Aethusa cynapium, Daucus carota, Equisetum arvense, Celosia argentea, Cynotis species such as Cynotis axilaris, Parthenium species such as Parthenium hysterophorus, Physalis species such as Physalis minima, Digera species such as Digera arvensis, and the like. [0046] The herbicidal composition comprising imazethapyr and cycloxydim is particular suitable for controlling monocotyledonous weeds, preferably Echinochloa spec., Leptochloa spec., Orysa spec., Phalaris spec., Brachiaria spec., more preferably Echinochloa spec. such as Echinochloa colonum, Leptochloa spec. such as Leptochloa fusca and Orysa spec. such as Orysa sativa. [0047] The herbicidal comprising imazethapyr, imazamox and cycloxydim is particular suitable for Echinochloa spec., Leptochloa spec., Orysa spec., Phalaris spec., Brachiaria spec., more preferably Phalaris spec. such as Phalaris canariensis, Echinochloa spec. such as Echinochloa colonum, Phalaris spec. such as Phalaris canariensis, Orysa spec. such as Orysa sativa and Brachiaria spec. such as Brachiaria plantaginea. [0048] The combination of imazethapyr and cycloxydim is in particular useful against annual and perennial grasses and broad-leaved weeds in post-emergence application. [0049] The compositions of the present invention are also suitable for controlling a large number of annual and perennial sedge weeds including cyperus species such as purple nutsedge ( Cyperus rotundus L.), yellow nutsedge ( Cyperus esculentus L.), hime-kugu ( Cyperus brevifolius H.), sedge weed ( Cyperus microiria Steud), rice flatsedge ( Cyperus iria L.), and the like. [0050] The compositions according to the present invention are suitable for combat-ing/controlling common harmful plants in useful plants (i.e. in crops). The compositions of the present invention are generally suitable for combating/controlling undesired vegetation in Legumes (Fabaceae), including e.g. soybeans ( Glycine max ), peanuts ( Arachis hypogaea ) and pulse crops such as peas including Pisum sativum, pigeon pea and cowpea, beans including broad beans ( Vicia faba ), Vigna spp., and Phaseolus spp. and lentils ( lens culinaris var.). crops which are tolerant to the action of acetohydroxyacid synthase inhibiting herbicides, such as for example Clearfield® wheat, Clearfield® barley, Clearfield® corn, Clearfield® lentil, Clearfield® oilseed rape or canola, Clearfield® rice, Cultivance® soybean and/or Clearfield® sunflower. [0053] The compositions of the present invention are in particular suitable for combat-ing/controlling undesired vegetation in soybean, peanut, pea, bean, lentil, green gram, black gram, cluster bean, fenugreek, other pulse or leguminous crops, or crops which are tolerant to the action of acetohydroxyacid synthase inhibiting herbicides, such as for example Clearfield® wheat, Clearfield® barley, Clearfield® corn, Clearfield® lentil, Clearfield® oilseed rape or canola, Clearfield® rice, Cultivance® soybean and/or Clearfield® sunflower. [0054] The compositions of the present invention are most suitable for combating/controlling undesired vegetation in soybean, peanut, pea, bean, lentil, green gram, black gram, cluster bean, fenugreek, other pulse or leguminous crops, preferably soybean ( Glycine max ). [0055] If not stated otherwise, the compositions of the invention are suitable for application in any variety of the aforementioned crop plants. [0056] The compositions according to the invention can also be used in crop plants which are tolerant to one or more herbicides owing to genetic engineering or breeding, which are resistant to one or more pathogens such as plant pathogenous fungi to genetic engineering or breeding, or which are resistant to attack by insects owing to genetic engineering or breeding. Suitable are for example pulse or leguminous crop plants, preferably soybean, peanut, pea, bean, green gram, black gram, cluster bean, fenugreek, or lentil which are tolerant to herbicidal AHAS inhibitors, such as, for example imazethapyr, or pulse or leguminous crop plants, preferably soybean, peanut, pea, bean, green gram, black gram, cluster bean, fenugreek, or lentil which, owing to introduction of the gene for Bt toxin by genetic modification, are resistant to attack by certain insects. Most suitable are soybeans which are tolerant to herbicidal AHAS inhibitors, such as, for example imazethapyr or soybean plants which, owing to introduction of the gene for Bt toxin by genetic modification, are resistant to attack by certain insects. [0057] The compositions of the present invention can be applied in conventional manner by using techniques as skilled person is familiar with. Suitable techniques include spraying, atomizing, dusting, spreading or watering. The type of application depends on the intended purpose in a well known manner; in any case, they should ensure the finest possible distribution of the active ingredients according to the invention. [0058] The compositions can be applied pre- or post-emergence, i.e. before, during and/or after emergence of the undesirable plants. Preferably, the compositions are applied post-emergence, in particular after the emergence of both the crop and the undesirable plants. When the compositions are used in crops, they can be applied after seeding and before or after the emergence of the crop plants. The compositions invention can, however, also be applied prior to seeding of the crop plants. [0059] In any case imazethapyr and cycloxydim and optionally imazamox can be applied simultaneously or in succession. [0060] Also imazethapyr, cycloxydim, bentazone and optionally imazamox can be applied simultaneously or in succession. [0061] The compositions are applied to the plants mainly by spraying, in particular foliar spraying. Application can be carried out by customary spraying techniques using, for example, water as carrier and spray liquor rates of from about 10 to 2000 l/ha or 50 to 1000 l/ha (for example from 100 to 500 l/ha). Application of the herbicidal compositions by the low-volume and the ultra-low-volume method is possible, as is their application in the form of microgranules. [0062] In the case of a post-emergence treatment of the plants, the herbicidal mixtures or compositions according to the invention are preferably applied by foliar application. Application may be effected, for example, by usual spraying techniques with water as the carrier, using amounts of spray mixture of approx. 20 to 1000 l/ha. [0063] The required application rate of the composition of the pure active compounds, i.e. of imazethapyr, cycloxydim and optionally imazamox and/or safener (imazethapyr, cycloxydim and bentazone and optionally imazamox and/or safener respectively) depends on the density of the undesired vegetation, on the development stage of the plants, on the climatic conditions of the location where the composition is used and on the application method. In general, the application rate of the composition (total amount of imazethapyr, cycloxydim and optional further actives) is from 1 to 2,000 g/ha, preferably from 5 to 500 g/ha of active substance. [0064] The required application rates of imazethapyr are generally in the range from 1 g/ha to 200 g/ha and preferably in the range from 5 g/ha to 150 g/ha or from 10 g/ha to 100 g/ha of active substance. [0065] The required application rates of cycloxydim are generally in the range from 1 g/ha to 200 g/ha and preferably in the range from 2 g/ha to 150 g/ha or from 5 g/ha to 100 g/ha of active substance. [0066] The required application rates of imazamox are generally in the range from 1 g/ha to 100 g/ha and preferably in the range from 5 g/ha to 75 g/ha or from 10 g/ha to 50 g/ha of active substance. [0067] The imazethapyr/cycloxydim combination is usually in a weight ratio range of 100:1 to 1:100, preferably of 10:1 to 1:10, or 1:3 to 3:1, more preferably 1:2 to 2:1. [0068] In those combinations which contain imazamox in addition to imazethapyr and cycloxydim, the weight ratio of imazethapyr to imazamox is usually in the range of 10:1 to 1:10, preferably of 5:1 to 1:5, more preferably 2:1 to 1:2. [0069] In those combinations which contain imazamox in addition to imazethapyr and cycloxydim, the weight ratio of the total amount of imazethapyr and imazamox to cycloxydim is 1:20 to 20:1 or 1:10 to 10:1, more preferably 1:2 to 2:1. [0070] In those combinations which contain imazethapyr, cycloxydim and bentazone, the weight ration of imazethapyr to bentazone is from ranges from 1:500 to 10:1 or 1:100 to 1:1 or 1:20 to 1:3 [0071] In those combinations which contain imazethapyr, cycloxydim and bentazone, the weight ration of cycloxydim to bentazone is from ranges from 1:500 to 10:1 or 1:100 to 1:1 or 1:20 to 1:3. [0072] In those combinations which contain imazethapyr, cycloxydim, bentazone and imazamox, the weight ratio of the total amount of imazethapyr and imazamox to bentazone is 10:1 to 1:500 or 1:1 to 1:100 or 1:3 to 1:20. [0073] The required application rates of the safener, if applied, are generally in the range from 1 g/ha to 2,000 g/ha and preferably in the range from 2 g/ha to 2,000 g/ha or from 5 g/ha to 2,000 g/ha of active substance. Preferably no safener or virtually no safener is applied. [0074] The compositions of these embodiments are particularly suitable for controlling mono- and dicotyledonous weeds and sedge weeds, in particular Acalypha spec., Physalis spec., Digera spec., Aegilops Cylindrica, Agropyron repens, Alopecurus myosuroides, Avena fatua, Brassica spec., Brachiaria spec., Bromus spec., Echinochloa spec. such as for example Echinochloa colonum, Lolium spec., Phalaris spec., red rice, Setaria spec., Leptochloa spec. such as Leptochloa fusca, Orysa spec. such as Orysa sativa, Sorghum spec., Abuthilon theoprasti, Amarantus spec., Brassica kaber, Capsella bursa - pastoris, Chenopodium spec., Cyperus spec., Euphorbia spec. such as for example Euphorbia hirta or Euphorbia geniculata, Geranium sepc., Ipomoea spec., Polygonum spec., Raphanus raphanistrum, Sinapis arevensis, Sysimbrium spec., Thlaspi arvense, Rottboellia cochinchinensis, Dinebra spec., Digitaria sanguinalis, Eleusine indica, Saccharum spontaneum, Cynodon dactylon, Commelina benghalensis, Commelina communis, Parthenium spec., Celosia argentea, Cynotis spp and Xanthium strumarium. [0075] The compositions of these embodiments are particularly suitable for controlling monocotyledonous weeds, preferably Echinochloa spec., Leptochloa spec., Orysa spec., Phalaris spec., Brachiaria spec. [0076] The compositions of these embodiments are in particular suitable for combating undesired vegetation in soybean, peanut, pea, bean, lentil, green gram, black gram, cluster bean, fenugreek, other pulse or leguminous crops, or crops which are tolerant to the action of acetohydroxyacid synthase inhibiting herbicides, such as for example Clearfield® wheat, Clearfield® barley, Clearfield® corn, Clearfield® lentil, Clearfield® oilseed rape or canola, Clearfield® rice, Cultivance® soybean and/or Clearfield® sunflower. The compositions of these embodiments are most suitable for combating undesired vegetation in soybean, peanut, pea, bean, lentil, green gram, black gram, cluster bean, fenugreek, other pulse or leguminous crops, preferably soybean. [0077] If not stated otherwise, the compositions of this embodiment are suitable for application in any variety of the aforementioned crop plants. [0078] In particular, the compositions of these embodiments are suitable for application in soybean, peanut, pea, bean, lentil, green gram, black gram, cluster bean, fenugreek, other pulse or leguminous crops, preferably soybean. [0079] The present invention also relates to formulations of the compositions according to the present invention. The formulations contain, besides the composition, at least one organic or inorganic carrier material. The formulations may also contain, if desired, one or more surfactants and, if desired, one or more further auxiliaries customary for crop protection compositions. [0080] The formulation may be in the form of a single package formulation containing both imazethapyr and cycloxydim and optionally imazamox together with liquid and/or solid carrier materials, and, if desired, one or more surfactants and, if desired, one or more further auxiliaries customary for crop protection compositions. The formulation may be in the form of a two package formulation, wherein one package contains a formulation of imazethapyr and optionally imazamox while the other package contains a formulation of cycloxydim and wherein both formulations contain at least one carrier material, if desired, one or more surfactants and, if desired, one or more further auxiliaries customary for crop protection compositions. In the case of two package formulations the formulation containing imazethapyr and optionally imazamox and the formulation containing cycloxydim are mixed prior to application. In case the imazethapyr and optionally imazamox itself is a two package formulation the composition is in the form of a three-pack formulation. Preferably the mixing is performed as a tank mix, i.e. the formulations are mixed immediately prior or upon dilution with water. [0081] The formulation may be in the form of a single package formulation containing imazethapyr, cycloxydim and bentazone and optionally imazamox together with liquid and/or solid carrier materials, and, if desired, one or more surfactants and, if desired, one or more further auxiliaries customary for crop protection compositions. [0082] The formulation may be in form of a two package formulation or three package formulation, wherein one package may contain more than one herbicide and wherein each package may contain at least one carrier material, if desired, one or more surfactants and, if desired, one or more further auxiliaries customary for crop protection compositions. The formulation may be in form or a three package formulation or four package formulation, wherein each package may contain one herbicide and wherein each package may contain at least one carrier material, if desired, one or more surfactants and, if desired, one or more further auxiliaries customary for crop protection compositions. The different packages are mixed prior to application. Preferably the mixing is performed as a tank mix, i.e. the formulations are mixed immediately prior or upon dilution with water. [0083] In the formulation of the present invention the active ingredients, i.e. imazethapyr, cycloxydim, optionally imazamox and optional further actives are present in suspended, emulsified or dissolved form. The formulation according to the invention can be in the form of aqueous solutions, powders, suspensions, also highly-concentrated aqueous, oily or other suspensions or dispersions, aqueous emulsions, aqueous microemulsions, aqueous suspoemulsions, oil dispersions, pastes, dusts, materials for spreading or granules. [0084] Depending on the formulation type, they comprise one or more liquid or solid carriers, if appropriate surfactants (such as dispersants, protective colloids, emulsifiers, wetting agents and tackifiers), and if appropriate further auxiliaries which are customary for formulating crop protection products. The person skilled in the art is sufficiently familiar with the recipes for such formulations. Further auxiliaries include e.g. organic and inorganic thickeners, bactericides, antifreeze agents, antifoams, colorants and, for seed formulations, adhesives. [0085] Suitable carriers include liquid and solid carriers. Liquid carriers include e.g. non-aqueous solvents such as cyclic and aromatic hydrocarbons, e.g. paraffins, tetrahydronaphthalene, alkylated naphthalenes and their derivatives, alkylated benzenes and their derivatives, alcohols such as methanol, ethanol, propanol, butanol and cyclohexanol, ketones such as cyclohexanone, strongly polar solvents, e.g. amines such as N-methylpyrrolidone, and water as well as mixtures thereof. Solid carriers include e.g. mineral earths such as silicas, silica gels, silicates, talc, kaolin, limestone, lime, chalk, bole, loess, clay, dolomite, diatomaceous earth, calcium sulfate, magnesium sulfate, magnesium oxide, ground synthetic materials, fertilizers such as ammonium sulfate, ammonium phosphate, ammonium nitrate, ureas, and products of vegetable origin such as cereal meal, tree bark meal, wood meal and nutshell meal, cellulose powders, or other solid carriers. [0086] Suitable surfactants (adjuvants, wetting agents, tackifiers, dispersants and also emulsifiers) are the alkali metal salts, alkaline earth metal salts and ammonium salts of aromatic sulfonic acids, for example lignosulfonic acids (e.g. Borrespers™-types, Borregaard), phenolsulfonic acids, naphthalenesulfonic acids (Morwet types, Akzo Nobel) and dibutylnaphthalenesulfonic acid (Nekal® types, BASF SE), and of fatty acids, alkyl- and alkylarylsulfonates, alkyl sulfates, lauryl ether sulfates and fatty alcohol sulfates, and salts of sulfated hexa-, hepta- and octadecanols, and also of fatty alcohol glycol ethers, condensates of sulfonated naphthalene and its derivatives with formaldehyde, condensates of naphthalene or of the naphthalenesulfonic acids with phenol and formaldehyde, polyoxyethylene octylphenol ether, ethoxylated isooctyl-, octyl- or nonylphenol, alkylphenyl or tributylphenyl polyglycol ether, alkylaryl polyether alcohols, isotridecyl alcohol, fatty alcohol/ethylene oxide condensates, ethoxylated castor oil, polyoxyethylene alkyl ethers or polyoxypropylene alkyl ethers, lauryl alcohol polyglycol ether acetate, sorbitol esters, lignosulfite waste liquors and proteins, denaturated proteins, polysaccharides (e.g. methylcellulose), hydrophobically modified starches, polyvinyl alcohol (Mowiol® types Clariant), polycarboxylates (BASF SE, Sokalan® types), polyalkoxylates, polyvinylamine (BASF SE, Lupamine® types), polyethyleneimine (BASF SE, Lupasol® types), polyvinylpyrrolidone and copolymers thereof. [0087] Examples of thickeners (i.e. compounds which impart to the formulation modified flow properties, i.e. high viscosity in the state of rest and low viscosity in motion) are polysaccharides, such as xanthan gum (Kelzan® from Kelco), Rhodopol® 23 (Rhone Poulenc) or Veegum® (from R.T. Vanderbilt), and also organic and inorganic sheet minerals, such as Attaclay® (from Engelhardt). [0088] Examples of antifoams are silicone emulsions (such as, for example, Silikon®-SRE, Wacker or Rhodorsil® from Rhodia), long-chain alcohols, fatty acids, salts of fatty acids, organofluorine compounds and mixtures thereof. [0089] Bactericides can be added for stabilizing the aqueous herbicidal formulations. Examples of bactericides are bactericides based on diclorophen and benzyl alcohol hemiformal (Proxel® from ICI or Acticide® RS from Thor Chemie and Kathon® MK from Rohm & Haas), and also isothiazolinone derivates, such as alkylisothiazolinones and benzisothiazolinones (Acticide® MBS from Thor Chemie). [0090] Examples of antifreeze agents are ethylene glycol, propylene glycol, urea or glycerol. [0091] Examples of colorants are both sparingly water-soluble pigments and water-soluble dyes. [0092] Examples which may be mentioned are the dyes known under the names Rhodamin B, C.I. Pigment Red 112 and C.I. Solvent Red 1, and also pigment blue 15:4, pigment blue 15:3, pigment blue 15:2, pigment blue 15:1, pigment blue 80, pigment yellow 1, pigment yellow 13, pigment red 112, pigment red 48:2, pigment red 48:1, pigment red 57:1, pigment red 53:1, pigment orange 43, pigment orange 34, pigment orange 5, pigment green 36, pigment green 7, pigment white 6, pigment brown 25, basic violet 10, basic violet 49, acid red 51, acid red 52, acid red 14, acid blue 9, acid yellow 23, basic red 10, basic red 108. [0093] Examples of adhesives are polyvinylpyrrolidone, polyvinyl acetate, polyvinyl alcohol and tylose. [0094] To prepare emulsions, pastes or oil dispersions, the active the components, as such or dissolved in an oil or solvent, can be homogenized in water by means of wetting agent, tackifier, dispersant or emulsifier. Alternatively, it is possible to prepare concentrates consisting of active substance, wetting agent, tackifier, dispersant or emulsifier and, if desired, solvent or oil, and these concentrates are suitable for dilution with water. [0095] Powders, materials for spreading and dusts can be prepared by mixing or concomitant grinding of the active the components imazethapyr and cycloxydim and optionally imazamox and optionally safener with a solid carrier. [0096] Granules, e.g. coated granules, impregnated granules and homogeneous granules, can be prepared by binding the active ingredients to solid carriers. [0097] The formulations of the invention comprise a herbicidally effective amount of the composition of the present invention. The concentrations of the active ingredients in the formulations can be varied within wide ranges. In general, the formulations comprise from 1 to 98% by weight, preferably 10 to 60% by weight, of active ingredients (sum of imazethapyr and cycloxydim and optionally imazamox and/or further acitves; sum of imazethapyr, cycloxydim and bentazone and optionally imazamox and/or further actives). The active ingredients are employed in a purity of from 90% to 100%, preferably 95% to 100% (according to NMR spectrum). [0098] The active compounds imazethapyr and cycloxydim and optionally imazamox as well as the compositions according to the invention can, for example, be formulated as follows: 1. Products for Dilution with Water [0099] A Water-Soluble Concentrates [0100] 10 parts by weight of active compound (or composition) are dissolved in 90 parts by weight of water or a water-soluble solvent. As an alternative, wetters or other adjuvants are added. The active compound dissolves upon dilution with water. This gives a formulation with an active compound content of 10% by weight. [0101] B Dispersible Concentrates [0102] 20 parts by weight of active compound (or composition) are dissolved in 70 parts by weight of cyclohexanone with addition of 10 parts by weight of a dispersant, for example polyvinylpyrrolidone. Dilution with water gives a dispersion. The active compound content is 20% by weight. [0103] C Emulsifiable Concentrates [0104] 15 parts by weight of active compound (or composition) are dissolved in 75 parts by weight of an organic solvent (eg. alkylaromatics) with addition of calcium dodecyl-benzenesulfonate and castor oil ethoxylate (in each case 5 parts by weight). Dilution with water gives an emulsion. The formulation has an active compound content of 15% by weight. [0105] D Emulsions [0106] 25 parts by weight of active compound (or composition) are dissolved in 35 parts by weight of an organic solvent (eg. alkylaromatics) with addition of calcium dodecyl-benzenesulfonate and castor oil ethoxylate (in each case 5 parts by weight). This mixture is introduced into 30 parts by weight of water by means of an emulsifier (Ultraturrax) and made into a homogeneous emulsion. Dilution with water gives an emulsion. The formulation has an active compound content of 25% by weight. [0107] E Suspensions [0108] In an agitated ball mill, 20 parts by weight of active compound (or composition) are comminuted with addition of 10 parts by weight of dispersants and wetters and 70 parts by weight of water or an organic solvent to give a fine active compound suspension. Dilution with water gives a stable suspension of the active compound. The active compound content in the formulation is 20% by weight. [0109] F Water-Dispersible Granules and Water-Soluble Granules [0110] 50 parts by weight of active compound (or composition) are ground finely with addition of 50 parts by weight of dispersants and wetters and made into water-dispersible or water-soluble granules by means of technical appliances (for example extrusion, spray tower, fluidized bed). Dilution with water gives a stable dispersion or solution of the active compound. The formulation has an active compound content of 50% by weight. [0111] G Water-Dispersible Powders and Water-Soluble Powders [0112] 75 parts by weight of active compound (or composition) are ground in a rotor-stator mill with addition of 25 parts by weight of dispersants, wetters and silica gel. Dilution with water gives a stable dispersion or solution of the active compound. The active compound content of the formulation is 75% by weight. [0113] H Gel Formulations [0114] In a ball mill, 20 parts by weight of active compound (or composition), 10 parts by weight of dispersant, 1 part by weight of gelling agent and 70 parts by weight of water or of an organic solvent are mixed to give a fine suspension. Dilution with water gives a stable suspension with active compound content of 20% by weight. [0115] 2. Products to be Applied Undiluted [0116] I Dusts [0117] 5 parts by weight of active compound (or composition) are ground finely and mixed intimately with 95 parts by weight of finely divided kaolin. This gives a dusting powder with an active compound content of 5% by weight. [0118] J Granules (GR, FG, GG, MG) [0119] 0.5 parts by weight of active compound (or composition) are ground finely and associated with 99.5 parts by weight of carriers. Current methods here are extrusion, spray-drying or the fluidized bed. This gives granules to be applied undiluted with an active compound content of 0.5% by weight. [0120] K ULV Solutions (UL) [0121] 10 parts by weight of active compound (or composition) are dissolved in 90 parts by weight of an organic solvent, for example xylene. This gives a product to be applied undiluted with an active compound content of 10% by weight. [0122] Aqueous use forms can be prepared from emulsion concentrates, suspensions, pastes, wettable powders or water-dispersible granules by adding water. [0123] It may furthermore be beneficial to apply the compositions of the invention alone or in combination with other herbicides, or else in the form of a mixture with other crop protection agents, for example together with agents for controlling pests or phytopathogenic fungi or bacteria. Also of interest is the miscibility with mineral salt solutions, which are employed for treating nutritional and trace element deficiencies. Other additives such as non-phytotoxic oils and oil concentrates may also be added. [0124] Synergism can be described as an interaction where the combined effect of two or more compounds is greater than the sum of the individual effects of each of the compounds. The presence of a synergistic effect in terms of percent control, between two mixing partners (X and Y) can be calculated using the Colby equation (Colby, S. R., 1967, Calculating Synergistic and Antagonistic Responses in Herbicide Combinations, Weeds, 15, 21-22): [0000] E = X + Y - XY 100 [0125] When the observed combined control effect is greater than the expected (calculated) combined control effect (E), then the combined effect is synergistic. [0126] The following tests demonstrate the control efficacy of compounds, mixtures or compositions of this invention on specific weeds. However, the weed control afforded by the compounds, mixtures or compositions is not limited to these species. The analysis of synergism or antagonism between the mixtures or compositions was determined using Colby's equation. Analogously, the Colby's equation can be used to determine synergism of 3-way and higher mixtures. EXAMPLES [0127] Products: [0128] Imazethapyr—70% WG [0129] Imazamox—70% WG [0130] (Imazethapyr 35+Imazamox 35)—70% WG—Ready mix [0131] Cycloxydim—200 g/lit EC [0132] Weeds in the Study [0000] EPPO Code Scientific Name ECHCO Echinochloa colonum LEFFA Leptochloa fusca ORYSA Orysa sativa PHACA Phalaris canariensis BRADE Brachiaria deflexa BRAPL Brachiaria plantaginea Example 1 Post Emergence Treatment by the Mixture of Imazethapyr+Cycloxydim [0133] [0000] Herbicidal activity against Application rate in g ai/ha ECHCO LEFFA ORYSA Imazethapyr Cycloxydim Found Calculated Found Calculated Found Calculated 7.5 — 85 — 55 — 50 — — 6 50 — 40 — 20 — 7.5 6 95 93 85 73 80 60 Example 2 Post Emergence Treatment by the Mixture of (Imazethapyr+Imazamox)+Cycloxydim [0134] [0000] Application rate in g ai/ha Herbicidal activity against (Imazethapyr + PHACA ECHCO Imazamox) Cycloxydim Found Calculated Found Calculated 30 — 50 — 98 — — 3 0 — 0 — 30 3 85 50 100 98 [0000] Application rate in g ai/ha Herbicidal activity against (Imazethapyr + PHACA ORYSA Imazamox) Cycloxydim Found Calculated Found Calculated 15 — 85 — 90 — — 6 60 — 20 — 15 6 95 94 95 92 [0000] Application rate in g ai/ha Herbicidal activity against (Imazethapyr + PHACA BRAPL Imazamox) Cycloxydim Found Calculated Found Calculated 15 — 50 — 80 — — 3 0 — 0 — 15 3 75 50 85 80 [0135] The invention further refers to the following embodiments: [0136] 1. Herbicidal compositions comprising imazethapyr, or an agriculturally acceptable salt thereof and cycloxydim, or an agriculturally acceptable salt thereof and optionally imazamox, or an agriculturally acceptable salt thereof. [0137] 2. Compositions according to embodiment 1, additionally containing a safener. [0138] 3. Compositions according to embodiment 1, containing no safener. [0139] 4. Compositions as according to any of the preceding embodiments, wherein the relative amount of imazethapyr to (optionally) imazamox and cycloxydim is in synergistically effective amounts. [0140] 5. Compositions as according to any of the preceding embodiments, wherein the relative amount of imazethapyr to cycloxydim is from 100:1 to 1:100, preferably 10:1 to 1:10. [0141] 6. The use of the compositions according to any of the preceding embodiments for controlling undesirable vegetation. [0142] 7. The use according to embodiment 6 for controlling undesirable vegetation in crop plants. [0143] 8. The use according to embodiment 7, wherein the crop plants are leguminous crops. [0144] 9. The use according to embodiment 8, wherein the crop plants are leguminous crops selected from soybean, green gram, black gram, peas, cluster beans, pulse crops, peanut. [0145] 10. The use according to embodiment 9, wherein the crop plant is soybean. [0146] 11. The use according to embodiment 7, wherein the crop plants are tolerant to acetohydroxyacid synthase inhibiting herbicides. [0147] 12. The use according to embodiment 11, wherein the crop plants tolerant to acetohydroxyacid synthase inhibiting herbicides are wheat, barley, canola, corn, lentils, oilseed rape, rice, soybean or sunflower. [0148] 13. A method for controlling undesirable vegetation, which comprises allowing a composition as according to embodiments 1 to 5 to act on plants to be controlled or their habitat. [0149] 14. A method for controlling undesired vegetation according to embodiment 13, which comprises applying the composition according to embodiments 1 to 5 before, during and/or after the emergence of the undesirable plants; the herbicides imazethapyr and cycloxydim and optionally imazamox being applied simultaneously or in succession. [0150] 15. An herbicide formulation comprising a composition according to embodiments 1 to 5 and at least one solid or liquid carrier.

The present invention relates to herbicidally active compositions, which comprise 5-ethyl-2-[(RS)-4-isopropyl-4-methyl-5-oxo-2-imidazolin-2-yl]nicotinic acid (common name: imazethapyr) and (5RS)-2-[(EZ)-1-(Ethoxyimino)butyl]-3-hydroxy-5-[(3RS)-thian-3-yl]cyclohex-2-en-1-one (common name: cycloxydim) and optionally 2-[(RS)-4-isopropyl-4-methyl-5-oxo-2-imidazolin-2-yl]-5-methoxymethylnicotinic acid (common name: imazamox).

TECHNICAL FIELD [0001] The invention relates to an apparatus and method for augmenting the three-dimensional position information obtained from the NAVSTAR satellite-based global positioning system (“GPS”) system. The apparatus and method demonstrates highest value when providing augmented position information to high-precision GPS-based guidance systems on blasthole drills that are routinely used in open pit (surface) mines. BACKGROUND [0002] Dead reckoning is a navigational technique which has been in use for centuries. Dead reckoning calculates the current position of an object based on a previous position of the object in view of the speed and direction travelled from the previous position. Disadvantageously, dead reckoning is subject to significant error, particularly when speed and direction are not measured accurately. [0003] The NAVSTAR (US government owned and operated) GPS constellation comprises a network of 27 Earth orbiting satellites. A complementary space-based network called GLONASS (Russian government owned and operated) consists of an additional 24 satellites. In order to determine the position of an object using GPS/GLONASS, a GPS, GLONASS or combined GPS/GLONASS receiver on the object must determine the location of at least four GPS/GLONASS satellites and the distance between the object and each of the at least four satellites. Disadvantageously, the GPS/GLONASS system cannot be used to calculate position when the GPS/GLONASS receiver does not receive signals from at least four GPS/GLONASS satellites. [0004] The introduction of high-precision global positioning systems (“HPGPS”) to the surface mining industry has resulted in significant improvements in productivity, and is expected to take an essential role as an enabling technology in future efforts to automate mining activities. In a standard system, GPS/GLONASS output is used directly for positioning. However, such systems can be impacted by physical obstacles that prevent the receipt of the satellite signals or as a result of sun spot activity that introduces noise into the signals thus causing them to become intermittently unavailable and/or making them less accurate in the course of normal operation. Therefore, an improved positioning solution that can operate under such poor GPS operational conditions is needed. The apparatus and method of the invention augments GPS with dead reckoning techniques when GPS signals are unavailable or inaccurate. SUMMARY OF THE INVENTION [0005] The invention relates to an augmented GPS (“aGPS”) apparatus and method, which alleviates the availability problem of a GPS receiver only (does not take into account a loss of GLONASS receiver signal) by combining it with dead reckoning techniques. The invention may be used in relation to a number of vehicle types (for example: tracked vehicles such as blasthole drills, excavators and bulldozers, or rubber tired vehicles such as haul-trucks and graders). The invention is particularly suited to blasthole drills. During operation, a blasthole drill typically travels for two minutes, stops and has its jacks lowered, drills for between thirty (30) and sixty (60) minutes, has its jacks retracted and travels an additional two (2) minutes where the process is repeated. In most cases, it is desirable for the blasthole drill to travel in a straight line for an extended period of time and distance. [0006] The augmented GPS of the invention introduces an intermediate step between the GPS receiver and the machine. Under normal conditions, the GPS output of the invention is identical to a GPS system. However, under degrading space-based GPS satellite conditions, the system instead estimates the motion of the machine based on local sensor measurements, and uses this to extrapolate the last known GPS position. This constructed position is output in place of the unavailable GPS position. This process of extrapolation continues until either the GPS situation improves, or the uncertainty in the constructed position exceeds a predefined maximum allowable value. The target precision of this system is to estimate the vehicle position within six (6) inches of its true value over a travelled distance of one hundred (100) feet (0.5%). [0007] In accordance with one aspect of the present invention, there is provided an augmented global positioning system (“aGPS”) for a vehicle comprising: (a) an aGPS computer; (b) a standard global positioning system (“GPS”) system operatively connected to said aGPS computer, the standard GPS comprising: (i) a high-precision GPS receiver; (ii) a navigation system; and (iii) a switch for alternating between a use of the standard GPS and the aGPS; and (c) a chorus subsystem operatively connected to the standard GPS, the chorus subsystem comprising: (i) a chorus data acquisition (DAQ) module; (ii) a gyroscope operatively connected to the chorus DAQ; and (iii) at least two rotation sensors operatively connected to the chorus DAQ. [0017] In accordance with another aspect of the present invention, there is provided a method for determining the position of a moving vehicle using augmenting global positioning system (“aGPS”), the method comprising the steps of: (a) calculating a first position of the vehicle using a global positioning system (“GPS”); (b) upon losing the GPS signal, measuring the movement of the vehicle and calculating the position of the vehicle using the last known position of the vehicle from the GPS combined with dead reckoning; (c) upon reacquiring a GPS signal, comparing the first position of the vehicle to the calculated position; and (d) correcting error in said calculated position. BRIEF DESCRIPTION OF THE FIGURES [0022] FIG. 1 is a schematic showing a first embodiment of the invention; [0023] FIG. 2 is a prospective side view showing an exemplary mounting location of a wheel sensor assembly; and [0024] FIG. 3 is a flow chart showing a method of augmenting GPS according to an embodiment of the invention. DETAILED DESCRIPTION [0025] Referring to FIG. 1 , a first embodiment of a high-availability global positioning system with local sensor augmentation ( 10 ) of the invention is shown. The system ( 10 ) is preferably used with blasthole drills, for example the Atlas Copco™ PV-271 blasthole drill. The system comprises a standard GPS system ( 20 ), an aGPS computer ( 30 ), and a chorus subsystem ( 40 ). [0026] The standard GPS system ( 20 ) comprises a dual-antenna high-precision GPS receiver ( 22 ), a navigation system ( 24 ), and a switch ( 26 ). The switch ( 26 ) allows the system to alternate between operation when a GPS signal is available, during which time the standard GPS system ( 20 ) is used, and when the GPS signal is not available, during which time the chorus subsystem ( 40 ) is used. [0027] The aGPS computer ( 30 ) acts as the processing unit for system ( 10 ), receiving sensor data as input and producing vehicle position information as output. [0028] The chorus subsystem ( 40 ) comprises a left rotation sensor ( 42 ), a right rotation sensor ( 44 ), and a chorus enclosure ( 50 ). Preferably, the left rotation sensor ( 42 ) and the right rotation sensor ( 44 ) are rotary encoders capable of measuring angular positions of the left and right wheels of the vehicle. The sensors ( 42 ) and ( 44 ) use a polarized magnet-sensor pair to sense the angular positions of the left and right drive motors, which directly drive the vehicle's crawler tracks. From sensors ( 42 ) and ( 44 ), the distance travelled by the vehicle is measured. The chorus enclosure ( 50 ) comprises a gyroscope ( 52 ) and a chorus data acquisition module ( 54 ). The gyroscope ( 52 ) obtains angular rate measurements about the vehicle's turning axis of rotation. For example, the gyroscope may be an ADIS16130 single-axis MEMS gyroscope produced by Analog Devices™. The chorus data acquisition module ( 54 ) comprises a supporting hardware unit which forwards sensor measurements from the left rotation sensor ( 42 ), right rotation sensor ( 44 ), and the gyroscope ( 52 ) to the aGPS computer ( 30 ). [0029] Referring to FIG. 2 , an exemplary mounting location of a wheel sensor assembly is shown. A magnetic wheel sensor assembly ( 60 ) is shown in association with a hydraulic propel motor ( 62 ) of a crawler track ( 64 ). The magnetic wheel sensor assembly ( 60 ) consists of a polarized magnet ( 66 ) and a nearby magnetic pickup sensor ( 68 ). The sensor is preferably a two-axis magnetometer (essentially a digital compass). The magnet ( 66 ) is rigidly attached to the wheel and rotates with it, thus the magnet's “north” rotates with the wheel. The sensor is able to sense the direction of this magnetic “north” as it rotates, thus providing an instantaneous angular position of the wheel. Alternatively, rotary encoders of any type capable of the required precision may be substituted for the magnet-based sensors. [0030] Referring to FIG. 3 , a method of the present invention is shown. During normal operation, a vehicle receives positional information from the standard GPS system ( 20 ) (Step 100 ). However, upon losing the GPS signal, the movement of the vehicle is measured and the new vehicle position is calculated in using chorus subsystem ( 40 ) (Step 200 ). This Step 200 comprises measuring the distance the vehicle has travelled using at least two wheel sensors. Step 200 further comprises measuring the direction the vehicle has travelled. Preferably, this is performed using at least one gyroscope ( 52 ). Step 200 may be repeated as necessary in response to intermittent GPS signals. Upon reacquiring a GPS signal, the first position of the vehicle is compared to the calculated position of Step 200 and any error in the calculated position is corrected (Step 300 ). Alternatively, the process may stop when the calculated position exceeds a predefined maximum allowable value (Step 400 ). [0031] The aGPS computer ( 30 ) contains a filter algorithm in order to maintain an optimal estimate of the position and orientation of the vehicle as it travels from point to point. The filter is an unscented Kalman filter (UKF)-based design incorporating wheel rotation sensors ( 42 , 44 ), a gyroscope ( 52 ), and a HPGPS ( 22 ) which is intermittently unavailable. Nomenclature [0032] In the following description, capital letters are used to denote quantities in an absolute “world” reference frame and lowercase letters to denote those in other reference frames. The global frame is a Cartesian frame predefined by the mine site and measured in metres. Mine site coordinates are specified in terms of a Northing (metres in the N direction), Easting (metres in the E direction), and an ellipsoidal height. For convenience, the “world frame” is a right-handed 3-D Cartesian frame comprising (X, Y, Z), where X is in the direction of the Easting, Y is in the direction of the Northing, and Z points upward and is related to the ellipsoidal height. [0033] A vehicle's local frame is defined similarly. It is a right-handed Cartesian frame comprising (x,y,z), where x is measured in the vehicle's “forward” direction, y is measured in the “leftward” direction, and z in the upward direction. The vehicle's frame is defined to be directly between the track midpoints, at ground level. Orientations are specified in terms of the coordinate axes. Rotations and orientations about the world frame's (X, Y, Z) axes are denoted Θ X , Θ Y , Θ Z respectively. Similarly, in the vehicle frame, Θ X , Θ Y , Θ Z are used for orientations. A hat (̂) is used to denote an estimated quantity. State [0034] Assuming the vehicle travels in a 2-D plane, only a subset of state variables are needed to achieve the desired accuracy. The state to be estimated is denoted q and consists of the global position and orientation of the vehicle's frame. It is represented as: [0000] q = [ X Y Θ Z ] . [0000] The state q has an associated 3×3 covariance matrix P. Initialization [0035] The above filter must be initialized using an absolute coordinate reference. Initialization can occur when two conditions are simultaneously met: 1. an RTK GPS fix is available. With this, the state variables X and Y can be initialized with the vehicle's current location in the absolute world coordinate frame; and 2. the vehicle is moving in a straight line, either forward or reverse. Since a single-antenna GPS receiver cannot measure its orientation, a heading is constructed based on consecutive GPS readings as detailed in the section entitled “Absolute Heading Estimate”, below. [0038] Since the RTK fix is not always available and since the vehicle spends most of its time stationary, it can take a long time for the above two conditions to be met under normal operating conditions. However, this can be remedied by making use of dual-antenna GPS hardware. The provision of a dual-antenna GPS hardware would remove condition “2” and allow the filter to initialize any time the RTK fix is available, regardless of the vehicle's motion. Absolute Heading Estimate [0039] While in theory, the GPS can reports its orientation via the HDT message, this is not a viable option likely due to the low speed of the drill. As an alternative, a heading can be constructed using the output GPS coordinates while the vehicle is moving. [0040] If the vehicle is moving, assuming two consecutive GPS coordinate readings (X 1 , Y 1 ) with uncertainty (σ X1 , σ Y1 ) and (X 2 , Y 2 ) with uncertainty (σ X2 , σ Y2 ), the heading Θ Z can be computed as [0000] β = Y 2 - Y 1  X 2 - X 1 ,  Θ Z = arctan  ( β ) . [0000] Using the standard error propagation formula, the uncertainties are [0000] σ β = 1 ( X 2 - X 1 ) 2  ( σ Y   2 2 + σ Y   1 2 ) + ( Y 2 - Y 1 ) 2 ( X 2 - X 1 ) 4  ( σ X   2 2 + σ X   1 2 ) ,  and σ Θ Z = σ β 1 + β 2 . [0000] Since the above process implicitly assumes that the vehicle is moving in a straight line (i.e. {dot over (Θ)} Z =0), an addition error component, σ m , is defined to account for error due to movement during the measurement. This additional error can be expressed as: [0000] σ m = d R - d L W , [0000] where d R and d L are the differential distances moved by the tracks during the measurement interval, and W is the distance between the tracks. If the vehicle is actually moving in a straight line, then d R ≈d L and σ m ≈0. Thus, σ Θ Z is defined to be: [0000] σ Θ Z = σ β 1 + β 2 + σ m . [0041] A number of conditions on the input data are enforced before applying the above procedure to construct a heading estimate. If any of these conditions fail, no computed heading is available. The conditions are: [0042] 1. Both GPS data points (X 1 , Y 1 ) and (X 2 , Y 2 ) must have RTK precision. [0043] 2. There is a minimum distance between the GPS data points. The distance d is computed using the formula: [0000] d =√{square root over (( X 2 −X 1 ) 2 +( Y 2 −Y 1 ) 2 )} [0044] The threshold used is d min =0.1 m. Thus, this condition is met if d≧d min . [0045] 3. The track speed of the left and right tracks must be similar. This confirms that the drill is travelling in a straight line, either forward or backward. The distance travelled by each track during the interval between data points denoted dr L and dr R are computed using the difference sin angular values with a constant found by calibration. The absolute value of their difference Δdr is then compared against a threshold value Δdr max . Absolute Position Estimate [0046] The absolute position X, Y, Z is obtained directly from the HPGPS' PTNL, PJK message. Since the GPS is not located at the defined machine origin, the reported values must be transformed into the machine frame using the most recent estimate of Θ Z . Defining the offset of the GPS antenna in the vehicle's frame as (x GPS , y GPS ), the absolute position of the GPS is [0000] [ X GPS Y GPS ] = [ X Y ] + [ cos   Θ Z - sin   Θ Z sin   Θ Z cos   Θ Z ]  [ x GPS y GPS ] ( 1 ) [0047] The corresponding covariance P is obtained directly from the GST message. This formulation can be extended to the full 3D case later if necessary. Vehicle Kinematic Model [0048] A tracked vehicle is modelled as a differential-drive vehicle with two wheels separated by a distance W. Using measurements from the wheel encoders and an experimentally-determined calibration constant, the differential distances each track has moved since the last step can be measured. For the right and left tracks respectively, these are Δr R and Δr L . The updated equation is [0000] q k + 1 =  q k + G s , k  u k =  [ x k y k θ z , k ] + [ 0.5  cos   θ z , k 0.5  cos   θ z , k 0.5  sin   θ z , k 0.5  sin   θ z , k 1 / W - 1 / W ]  [ Δ   r R Δ   r L ] .   ( 3 )   ( 2 ) State Update [0049] The filter's UKF-based estimation algorithm uses the familiar predict-update cycle to maintain its state estimate. [0050] The prediction (a-priori) step is always done and is based on dead reckoning measurements. The basic premise is to use the kinematic model, described above, in a UKF a-priori step with a modification incorporating both wheel encoders and the z-gyro as measurements for rotation. First, the rotation due to wheels Δθ w and the uncertainty σ w of the same is defined: [0000]   Δσ w = Δ   r L - Δ   r R W   σ w = ( F w  Δ   r R + C w ) 2 + ( F w  Δ   r L + C w ) 2 W , [0000] where F w and C w are constants. Next a simple condition is used to determine whether the vehicle is currently moving: [0000] Δ   r R + Δ   r L 2 ≥ D min [0000] where D min is a constant threshold. Depending on whether the condition (4) is true, one of the following is performed: [0051] 1. If (4) is true, the vehicle is moving. Thus a rotation measurement is obtained from the z-gyro: [0000] Δθ g =Tg z −b z [0000] σ g S g′ where g z is the current raw measurement from the gyroscope (in units of rad/s), T is the timestep, b z is the constant gyro bias (discussed below in step “2”), and S g is the constant uncertainty of the gyro measurements. Next the combined equivalent measurement and uncertainty as the uncertainty-weighted mean of the values from the wheels and the gyro is calculated: [0000] σ c = 1 1 σ w 2 + 1 σ g 2 .  Δθ c = σ c  ( Δσ w σ w 2 + Δθ g σ g 2 ) . With the uncertainties correctly adjusted, this scheme tends to trust the gyroscope measurements more while moving, and the wheel sensor measurements while moving slowly or stationary. [0054] 2. If (4) is false, we consider the vehicle to be stationary. In this case, the combined rotation values are those of the wheels alone: [0000] Δθ c =Δθ w [0000] σ c =σ w . [0000] Since the gyroscope's bias b z drifts over time, the (stationary) time can be used to estimate its current value. A (normal) Kalman filter is used to track both the gyroscope bias b z and its uncertainty σ b . The expression Δθ g −Δθ w represents a measurement of its current value, and incorporates it into b z using one step of the Kalman filter. This filter is effectively only an a-posteriori step. [0055] The a-priori step is then done using an unscented transformation, and incorporating Δθ c in place of Δθ w in the deconstructed model (2). [0056] The update step is done according to one of the cases below. [0057] 1. Case 1: RTK fix available and the conditions of “Absolute Heading Estimation” are fulfilled. In such a case, a heading is constructed as detailed in that section. The machine state q is transformed into the GPS frame using the inverse of Equation (1), and a full-state update is done in the GPS frame using the usual UKF update step with the recent GPS position and constructed heading. [0058] 2. Case 2: RTK fix available, but the conditions of “Absolute Heading Estimate” are not fulfilled. In such a case, the machine state q is transformed into the GPS frame using the inverse of Equation (1), and a partial-state update is done in the GPS frame using the usual UKF update step with the recent GPS (X, Y) position. [0059] 3. Case 3: RTK is not available, in which case, the update step is skipped. [0060] A person of skill in the art would recognize that the type, number, and position of said sensors and gyroscope may be varied according to the intended use. [0061] The scope of the claims should not be limited by the preferred embodiments set forth in the examples, but should be given the broadest interpretation consistent with the description as a whole.

The invention relates to an apparatus and method for augmenting the 3 dimensional position information obtained from the NAVSTAR satellite-based global positioning system (“GPS”) system. Such systems can be impacted by physical obstacles that prevent the receipt of the satellite signals or as a result of sun spot activity that introduces noise into the signals thus causing them to become intermittently unavailable and/or making them less accurate in the course of normal operation. Therefore, an improved positioning solution that can operate under such poor GPS operational conditions is needed. The apparatus and method of the invention augments GPS with dead reckoning techniques when GPS signals are unavailable or inaccurate. The apparatus and method of the invention demonstrates highest value when applied to blasthole drill positioning applications in open-pit mines.

BACKGROUND OF THE INVENTION 1. Field of the Invention The present inventions relate to dental hygiene and, more particularly, to toothbrushes. 2. Description of the Related Art To ensure proper oral care, dentists recommend that we brush our teeth more than once a day for at least two to three minutes each time. Despite this recommendation, the average adult person does not brush his or her teeth for two to three minutes. This problem is worse with children, who have notoriously short attention spans and often view brushing their teeth as a chore. Accordingly, there is a general need for a device that encourages people, especially children, to brush their teeth more often and for longer periods of time. See e.g., U.S. Patent Publication No. 2004-0143920, filed Jan. 24, 2003 and published on Jul. 29, 2004, the entirety of which is hereby incorporated by reference herein. SUMMARY OF THE INVENTION U.S. Patent Publication No. 2004-0143920 describes a toothbrush with a handle having a base, a body, and a head. The body can have a first section and a second section forming an oblique angle. A projector of a beam of light located within the handle. The toothbrush can have at least one bristle attached to the head. The toothbrush can have a grip attached to the base. An illumination circuit can be positioned within the handle and is operated by pressing an end of the handle. Pressing, the end of the handle can compress a spring which completes the illumination circuit, activating the projector of a light beam within the toothbrush. Applicant has recognized that some children may have trouble activating the illumination circuit. Accordingly, a need exists for an improved activation mechanism. Thus, one aspect of the present invention is a toothbrush comprising a handle having a first end and a second end and a head coupled to the first end of the handle, the head comprising a plurality of bristles. A pliant base is coupled to the second end of the handle. A light is positioned in the handle. The tooth brush also include a power source, a first contact member, a second contact member and a flexible member that extends around the first contact member and is coupled to the second contact member such that the second contact member contacts the first contact member as the pliant base is compressed or bent but does not contact the first member when the pliant base is in an unstressed condition. A control circuit configured such that contact between the first contact member and the second contact member completes a circuit and initiates illumination of the light for a set period of time. Another aspect of the present invention is a toothbrush comprising a handle having a first end and a second end and a head coupled to the first end of the handle, the head comprising a plurality of bristles. A pliant base is coupled to the second end of the handle. A light is positioned in the handle. The toothbrush also includes a power source, a first contact member, a second contact member and means for separating the first contact member from the second contact member flexible member and for permitting contact between the first and second contact member when the pliant base means is deflected in a direction transverse to the longitudinal axis. A control circuit configured such that contact between the first contact member and the second contact member completes a circuit and initiates illumination of the light for a set period of time. Yet another aspect of the present invention is a toothbrush comprising a handle having a first end and a second end and a head coupled to the first end of the handle, the head comprising a plurality of bristles. A pliant base is coupled to the second end of the handle. A light is positioned in the handle. The toothbrush also includes a power source, a first contact member that is generally stationary with respect to the power source and a second contact member that moves with respect to the first contact member as the actuation member is moved. A control circuit configured such that contact between the first contact member and the second contact member completes a circuit and initiates illumination of the light for a set period of time. The foregoing objects may also be achieved by a toothbrush having a handle having a base, a body, and a head. The body having a chamber therein. A projector of a beam of light within the chamber. The toothbrush having at least one bristle attached to the head. The toothbrush having a grip attached to the base. The foregoing objects may still further be achieved by a method of using a toothbrush. The method uses a toothbrush having a handle and an illumination circuit. The toothbrush handle having a base, a body, and a head. The toothbrush illumination circuit having a projector of a beam of light connected to a switch. The method including the step of gripping the toothbrush. The method further including the step of engaging the switch for completing the illumination circuit. The method still further including the step of activating a projector of a light beam within the toothbrush. The method still further including the step of utilizing the toothbrush while the projector of a light beam is activated. BRIEF DESCRIPTION OF THE DRAWINGS FIG. 1 is a perspective view of a prior art toothbrush. FIG. 2 is a front elevation view of the toothbrush of FIG. 1 showing the brush side of the toothbrush. FIG. 3 is a rear elevation view of the toothbrush of FIG. 1 showing the non-brush side of the toothbrush. FIG. 4 is a side view of the toothbrush of FIG. 1 showing the chamber preferentially placed near the non-brush side. FIG. 5 is a top view of the toothbrush with the grip removed exposing the toothbrush base and illumination circuit held within. FIG. 6 is an exploded top view of the toothbrush grip, base, and lower section of toothbrush handle showing the placement of the illumination circuit within the toothbrush handle. FIG. 7 is a sectional view of the grip of FIG. 6 . FIG. 8 is a sectional view of the base in FIG. 6 showing the indentations for containing the positive terminal conductors. FIG. 9 is a perspective view of the illumination circuit without the power supply. FIG. 10 is a bottom view of the illumination circuit. FIG. 11 is a top perspective view of the illumination circuit. FIG. 12 is a schematic drawing of the illumination circuit. FIG. 13 is a brush side perspective view of a toothbrush having certain features and advantages according to the present invention. FIG. 14 is top brush side exploded view of the tooth brush of FIG. 13 . FIG. 15 is a closer view of section 15 of FIG. 14 . FIG. 16 is another embodiment of an illumination circuit. FIG. 17 is another embodiment of an illumination circuit. FIG. 18 is a perspective view of an embodiment of a toothbrush with a front-mounted button. FIG. 19 is a perspective view of another embodiment of a toothbrush with a front-mounted button. FIG. 20 is a perspective view of another embodiment of a toothbrush with a front-mounted button. DETAILED DESCRIPTION OF THE PREFERRED EMBODIMENTS FIG. 1 illustrates a prior art illuminated toothbrush 10 , which comprises a handle 12 , an illumination circuit 14 , a brush 16 , and a grip 18 . As shown in FIG. 1 , the handle 12 comprises a base 20 , a body 21 , and a head 26 . The body 21 has a first section 22 and a second section 24 . The handle 12 can be formed of hard, clear plastic. In one arrangement, the handle 12 can be a colored plastic. In another arrangement, the handle 12 can be a translucent plastic. In yet another embodiment, the handle 12 may be fashioned out of a plastic incorporating metallic flake 55 . The toothbrush handle 12 can be formed through an injection molding process. In such an embodiment, plastic in a liquid form can be injected into a mold having two sections. Liquid plastic can be injected into the mold where it is then allowed to solidify. When the mold is opened it creates a handle having a brush side 28 and a non-brush side 30 . At the intersection of these two sides 28 , 30 can be a ridge 32 . The ridge 32 can be a surface characteristic resulting from the injection molding process. In the illustrated arrangement, the ridge 32 does not extend inside the handle 12 but exists on the surface. The injection molding process in constructing of the toothbrush handle 12 is conventional and does not form a part of the present invention. As seen in FIGS. 5 , 6 , and 8 , the handle 12 can include a handle base 20 . The base 20 can be generally cylindrical in shape and can have a circumferential groove or cavity 44 therein. The circumferential groove 44 can have a centerline. The base 20 has an outside surface 34 , an inside surface 36 , a first end 38 , and a second end 40 . In the illustrated arrangement, on the outside surface 34 of the base 20 are three annular rings 42 . The annular rings 42 can provide a ledge upon which the grip 18 holds. The inside surface 36 defines the cavity 44 . The inside surface 36 can have a first indentation 46 and a second indentation 48 . The cavity 44 can serve as a housing for the illumination circuit 14 . The first and second indentation 46 , 48 can serve as a guide for positioning the illumination circuit 14 within the base 20 . A chamber 50 can extend within the section 22 . The chamber 50 has a first end 52 and a second end 54 . The chamber 50 can be generally cylindrical in shape. The first end 52 can be rounded and can provide a transition between the chamber 50 filled full of air and the first section 22 which is of plastic. The air is inherently present as a result of assembly at a time after the handle was formed. Alternatively, the chamber may be filled full of a material in a process separate from the forming of the handle. The second end 54 of the chamber 50 can be open to the first end 38 of the base 20 . The first section chamber 50 can be in off-center alignment with the base 20 . As seen in FIG. 8 , the first section chamber 50 is in off-center alignment with the base 20 to allow the projector of a light beam or illuminating member 60 to emit a light beam 56 that travels through the first section 22 and to strike the interface 65 between the second section 24 and atmosphere. At this interface, a light beam 57 can be reflected towards the handle head 26 and a light beam 58 can be refracted towards the atmosphere. The first section chamber 50 can be also positioned in off-center alignment with the base 20 because in the first section 22 is ergometrically designed to accommodate a user's grip. In the ergometric design, the brush side 28 of the first section 22 is contoured and the non brush side 30 of the first section 22 is flat. In addition, the brush side 28 of the first section 22 arrives at a point of the second section 24 at a greater angle than the non brush side 30 . In other words, the illumination circuit 14 extends within the first section 22 substantially parallel to the center line of the base member 20 but the first section 22 brush side 28 angles toward the inner point where the first section 22 meets the second section 24 and the non brush side 30 portion of the first section 22 also angles toward the point where the first section 22 meets the second section 24 . Thus, for the first section chamber 50 to extend the furthest into the first section 22 of the handle 12 , the first section chamber 50 is preferably positioned closer to the non-brush side 30 of the first section 22 . The illumination circuit 14 can have an illuminating member or projector of a light beam 60 , a resistor 62 , a timing circuit 64 , and a power source 66 . These parts can be joined by the conductor 68 , which provides a support structure extending the illuminating member 60 a distance away from the timing circuit 64 . The negative terminal conductor 70 can be a spring which presses against the power source 66 , which in the illustrated arrangement comprises a series of batteries. The positive terminal conductor 72 can comprise a pair of prongs that extends away from the timing circuit 64 to embrace the power source 66 . The positive terminal conductor 72 can also be sized to stabilize the illumination circuit 14 within the base 20 as the positive terminal conductor 72 is sized to fit within the first indentation 46 and the second indentation 48 of the base 20 . The illumination member 60 in this embodiment is a light emitting diode (LED). In other embodiments, the illumination member 60 could be an incandescent light bulb. In still other embodiments, the illumination member 60 may be any other device known in the art that may provide illumination. The power source 66 in one arrangement can be micro-cell battery model number G3-ACNB. In the illustrated arrangement, three batteries are placed in series within the base 20 . The timing circuit 64 preferably can function to illuminate the illumination member 60 for approximately 60 seconds. The timing circuit 64 also preferably can serve to control the illumination member 60 to blink intermittently for the time period in which it is engaged. In some embodiments, the illumination member 60 may stay on continuously and/or illuminate for a longer or shorter period of time. In the illustrated embodiment, the circuit is activated by closing an electrical switch 74 to complete a circuit. The brush 16 can have a bristle 80 . The bristle 80 can have a first end 82 and a second end 84 . The bristle second end 84 can be embedded in the head 26 of the handle 12 . In one arrangement, the bristle 80 can be made of clear plastic material. Moreover, in such an arrangement, the brush 10 can be configured such that, when the bristle 80 is struck by light traveling from the illumination member 60 through the first section 22 and the second section 24 , a portion of the light striking the bristle 80 may reflect through the bristle 80 and extend out of the bristle 80 . The grip 18 can be made of a flexible material. In the illustrated arrangement, the grip 18 can also serve as a switch. For example, the grip 18 can have an extending piece or switch 74 of flexible material as seen in FIG. 7 . When the grip is pushed in the direction of arrow A in FIG. 4 , the piece 74 moves the positive terminal metal conductor 72 to contact the power source 66 . Alternatively, the piece 74 moves the power source 66 to contact the positive terminal metal conductor 72 . In doing so, the piece 74 moves illumination circuit 14 from an un-illuminated position to a illuminated position. The grip 18 can remain in place on the base 20 by engaging the annular rings 42 on the outside surface 34 of the base 20 . An adhesive 88 can be positioned between the first section 22 and the grip 18 to hold the grip 18 in close connection with the first section 22 . The grip 18 can be made of a flexible material. Alternatively, the grip 18 may be of a hard material but have a flexible portion that may be used to engage the positive terminal conductor 72 and press against the power source 66 . In operation, the illuminated toothbrush 10 is used by a user to indicate the duration of an amount of time. The user grips the toothbrush handle 12 in their hand with the bristle 80 surface with the bristle 80 against their teeth and engages the illumination circuit switch 74 . The illumination member 60 begins to blink intermittently in an on/off fashion. The illumination member 60 continues to blink for a period of approximately 60 seconds. The handle is designed to direct light to the user in multiple ways so that the user may be accurately apprised of brushing time. The frequency of blinking can remain constant, or vary in frequency. In some embodiments, the frequency can increase as the time approaches 60 seconds. In some embodiments, the frequency can remain constant through a first period of time, and increase in frequency in a second period of time. In one example, the frequency can remain constant for approximately 45 seconds; then increase for the remaining 15 seconds. In other embodiments, different time intervals can be used, such as, for example, two even periods of thirty seconds each. A light beam 59 travels from the illuminating device 14 through a first section chamber 50 . The light beam 59 strikes an interface 63 between the first section chamber and the first section and a light beam 61 is partially reflected off of the interface and a light beam 56 is refracted enters the first chamber. The light beam 56 travels through the first section 22 to strike upon an interface 65 between the second section 24 and atmosphere. A light beam 57 is reflected from interface 65 toward head 26 and a light beam 58 is refracted towards the atmosphere. The light beam 57 then strikes an interface 116 between the head 26 and bristle 80 where it is partially reflected and refracted. Alternatively, the toothbrush handle may have a metallic piece 55 or flake embedded in the hardened plastic. The angle of reflection upon the flake is equal to the angle of incidence upon the flake. These metallic pieces 55 can be glitter. In operation, light will strike these metallic pieces 55 at an angle of incidence and the reflected light beam 67 directed at an angle of reflection as seen in FIG. 1 . One disadvantage of the arrangement described above is that the mechanism for completing the activation of the illumination is mechanically inefficient and often requires a degree of strength and dexterity not possessed by children. FIGS. 13 , 14 , and 15 illustrate a modified embodiment of a toothbrush that advantageously addresses the aforementioned problem. Numerical reference to components is the same as in the previously described arrangement, except that a prime symbol (′) has been added to the reference. Where such references occur, it is to be understood that the components are the same or substantially similar to previously-described components. As can be seen the toothbrush can have an improved light generation mechanism 100 . The mechanism 100 can be disposed in the base 18 ′, as described above. In the illustrated embodiment, the mechanism 100 comprise a power portion 102 , an extension portion 104 , and an activation portion 106 . The illuminating member 60 ′ can be disposed at one end of one or more extension members 120 , which can form the extension portion and can extend toward the base 18 ′ and couple with the power segment 102 . The power segment 102 can comprise one or more power sources (e.g., batteries) 66 ′. The power sources 66 ′ can be positioned between a distal member 105 and a proximal member 109 , which in one embodiment can each comprise a circular disk-like plate. The power sources 66 ′ can be secured in the space between the distal and proximal members 105 , 109 and can be engaged by one or more generally rigid elongate members 107 . The elongate members 107 can extend along the longitudinal axis of the toothbrush 10 ′. For additional security, the power sources 66 ′ can be surrounded by a cylindrical member (not shown) that can be made of plastic or other similar material. In this manner, the power sources 66 ′ cannot be easily dislodged from the power segment and swallowed by children if the base 18 ′ is removed. The power sources 66 ′ can be in direct contact with each other or have an intervening electrical connection member (not shown). The power segment 102 can be coupled to an activation segment 106 . The activation segment 106 can have a spring member 108 . The spring member 108 can be a conical, as in the illustrated embodiment, or cylindrical, or any other shape appropriate for the interior or the base 18 ′. The spring member 108 can be composed of metal, though other materials can be used in other embodiments. Advantageously, an electrically-conducting material can be used. The spring member 108 can have an inward-extending protrusion 112 at one end. The protrusion 112 can be of any size or shape sufficient to extend towards the power sources 66 ′ without contacting them. The protrusion 112 can be composed of the same material as the spring member 108 , or can be composed of a different material, preferably an electrically-conductive material. In certain embodiments, the spring member 108 can be composed of a plastic and the protrusion 112 can be a metal. The spring member 108 can have an insulating or conducting coating. The power segment 102 can have a contact member 110 extending toward the protrusion 112 . The contact member 110 can be electrically-connected to the power sources 66 ′ and the illumination circuit 14 ′. The contact member 110 can be connected such that contact with the protrusion 112 activates the illumination circuit 14 ′. Additionally, when electrically-conducting materials are used for the spring member 108 , such as the metal in the illustrated embodiment, contact between the sides of the flexible member 108 and the contact member 110 can also activate the illumination circuit 14 ′. Thus, although the illustrated embodiment is shown in FIGS. 14 and 15 in an exploded view, when the toothbrush 10 ′ is assembled as in FIG. 13 , manipulation of the pliable base 18 ′ can cause deflection of the spring member 108 within. If the bottom of the base 18 ′ is pushed toward the power segment 102 , the protrusion 112 can touch the contact member 110 , causing illumination. Alternatively, if the base 18 ′ is deflected towards either side, the interior of the spring member 108 can touch the contact member 110 , also causing illumination. The illumination can be continuous or intermittent. Additionally, the intervals between illuminations during intermittent operation can be regular or have increasing or decreasing frequency. FIG. 16 illustrates another embodiment of an illumination circuit. Numerical reference to components is the same as in previously described arrangements, except that a double prime symbol (″) has been added to the reference. Where such references occur, it is to be understood that the components are the same or substantially similar to previously-described components. As in other embodiments, a contact member 110 ″ can be attached to a power segment 102 ″. The contact member 110 ″ can activate a circuit 14 ″, as illustrated in FIG. 12 . The embodiment of a mechanism 100 ″ depicted in FIG. 16 has a plurality of flexible members 212 which enclose the contact member 110 ″. The flexible members 212 can be composed of metal, an elastomer, or any of a variety of other materials which permit flexibility and have, or can support, an electrically-conductive surface. An inward-extending protrusion 112 ″ can be integrally formed with the flexible members 212 . Unlike the cone-shaped spring member depicted in the embodiment illustrated in FIG. 15 , the flexible members 212 can extend along a longitudinal axis of the mechanism 100 ″. In some embodiments, the flexible members 212 taper inwardly as they extend away from the contact member 110 ″. In other embodiments, the flexible members 212 do not taper, and maintain an approximately cylindrical shape. In yet other embodiments, the flexible members 212 can have other arrangements, including without limitation, a pyramidal prism, a rectangular prism, a cubic shape, or other geometrical shapes sized appropriately to surround the contact member 110 ″. FIG. 17 illustrates another embodiment of an illumination circuit 100 ′″. Numerical reference to components is the same as in previously described arrangements, except that a triple prime symbol (′″) has been added to the reference. Where such references occur, it is to be understood that the components are the same or substantially similar to previously-described components. The contact member 110 ′″ can be enclosed within a flexible mesh, such as a metal wire mesh 312 . The mesh 312 can have an interior contact surface which activates the illumination circuit 100 ′″, or can structurally support such a surface. The wire mesh 312 can have an inward-extending protrusion 112 ′″, as illustrated. In some embodiments, more than one protrusion is present on the interior of the component disposed around the contact member 110 ′″. These embodiments can include the use of spring members, flexible rods, flexible meshes, or any other contact surface or surface support configured to activate the illumination circuit. FIGS. 18-20 illustrate alternative embodiments of the toothbrush having a front-mounted activation mechanism for activating an illumination circuit 14 . The mechanism can comprise a variety of devices, some examples of which are illustrated and described below. FIG. 18 illustrates an embodiment of a toothbrush 410 having an illumination member 460 and an activation mechanism 468 . The mechanism 468 can comprise a contact port 470 and a button 472 . The button 472 can comprise a metallic mesh 474 that surrounds the contact post 470 , and activates an illumination circuit 14 , lighting the illumination member 460 , as described above. The mesh 474 can case the illumination circuit 14 to activate through contact with an electrically-conducting inner surface, or support an electrically-conducting surface which activates the circuit 14 . The mesh 474 can be replaced by a spring, flexible rods, or any other suitable device, as described above. FIG. 19 illustrates another embodiment of a toothbrush 510 having a front-mounted activation mechanism. The mechanism can comprise a push-button device 568 having a button 572 and a switch device 574 , as are well-known in the art. The push-button device 568 can cause the illumination member 560 to blink by activating an illumination circuit 14 . The switch device 574 can be activated by manipulation of the button 572 , whether the button 572 is flexible or a rigid connection to the switch device 574 . The push-button device 568 can activate the circuit 14 once manipulated, and future manipulations can be ignored by the circuit 14 until the timer has completed a cycle. This operation can occur in any embodiment described herein. FIG. 20 illustrates another embodiment of a toothbrush 610 , wherein an illumination member 660 is set to blink by an illumination circuit 14 . The circuit 14 can start a timed cycle upon receiving a signal from an activation device 668 . In the illustrated embodiment, the activation device 668 comprises a base 670 and two contact terminals 672 . The contact terminals 672 can activate the circuit 14 when electrical conduction occurs between the terminals 672 . In one embodiment, the circuit 14 and terminals 672 can be constructed to allow contact with human skin to both terminals 672 to cause conduction to occur, thereby activating the circuit 14 . In non-limiting examples, the palm of a human hand gripping the toothbrush can activate the circuit or, a finger or thumb pressed to touch both terminals 672 can activate the circuit 14 . Water disposed in continuous contact with both terminals 672 can also activate the circuit 14 . Although certain embodiments, features, and examples have been described herein, it will be understood by those skilled in the art that many aspects of the methods and devices shown and described in the present disclosure may be differently combined and/or modified to form still further embodiments. For example, any one component of the infusion sets shown and described above can be used alone or with other components without departing from the spirit of the present invention. Additionally, it will be recognized that the methods described herein may be practiced in different sequences, and/or with additional devices as desired. Such alternative embodiments and/or uses of the methods and devices described above and obvious modifications and equivalents thereof are intended to be included within the scope of the present invention. Thus, it is intended that the scope of the present invention should not be limited by the particular embodiments described above, but should be determined only by a fair reading of the claims that follow.

A toothbrush with a handle having a base, a body, and a head. The body having a first section and a second section forming an oblique angle. A projector of a beam of light located within the handle. The toothbrush having at least one bristle attached to the head. The toothbrush having a grip attached to the base. A method of using a toothbrush including the step of gripping the toothbrush. The method further including the step of engaging the projector of a beam of light. The method still further including the step of utilizing the toothbrush while the projector of a light beam is activated.