DreamBooth training example for Stable Diffusion XL (SDXL)

DreamBooth is a method to personalize text2image models like stable diffusion given just a few (3~5) images of a subject.

The train_dreambooth_lora_sdxl.py script shows how to implement the training procedure and adapt it for Stable Diffusion XL.

💡 Note: For now, we only allow DreamBooth fine-tuning of the SDXL UNet via LoRA. LoRA is a parameter-efficient fine-tuning technique introduced in LoRA: Low-Rank Adaptation of Large Language Models by Edward J. Hu, Yelong Shen, Phillip Wallis, Zeyuan Allen-Zhu, Yuanzhi Li, Shean Wang, Lu Wang, Weizhu Chen.

Running locally with PyTorch

Installing the dependencies

Before running the scripts, make sure to install the library's training dependencies:

Important

To make sure you can successfully run the latest versions of the example scripts, we highly recommend installing from source and keeping the install up to date as we update the example scripts frequently and install some example-specific requirements. To do this, execute the following steps in a new virtual environment:

git clone https://github.com/huggingface/diffusers
cd diffusers
pip install -e .

Then cd in the examples/dreambooth folder and run

pip install -r requirements_sdxl.txt

And initialize an 🤗Accelerate environment with:

accelerate config

Or for a default accelerate configuration without answering questions about your environment

accelerate config default

Or if your environment doesn't support an interactive shell (e.g., a notebook)

from accelerate.utils import write_basic_config
write_basic_config()

When running accelerate config, if we specify torch compile mode to True there can be dramatic speedups.

Dog toy example

Now let's get our dataset. For this example we will use some dog images: https://huggingface.co/datasets/diffusers/dog-example.

Let's first download it locally:

from huggingface_hub import snapshot_download

local_dir = "./dog"
snapshot_download(
    "diffusers/dog-example",
    local_dir=local_dir, repo_type="dataset",
    ignore_patterns=".gitattributes",
)

This will also allow us to push the trained LoRA parameters to the Hugging Face Hub platform.

Now, we can launch training using:

export MODEL_NAME="stabilityai/stable-diffusion-xl-base-1.0"
export INSTANCE_DIR="dog"
export OUTPUT_DIR="lora-trained-xl"

accelerate launch train_dreambooth_lora_sdxl.py \
  --pretrained_model_name_or_path=$MODEL_NAME  \
  --instance_data_dir=$INSTANCE_DIR \
  --output_dir=$OUTPUT_DIR \
  --mixed_precision="fp16" \
  --instance_prompt="a photo of sks dog" \
  --resolution=1024 \
  --train_batch_size=1 \
  --gradient_accumulation_steps=4 \
  --learning_rate=1e-4 \
  --report_to="wandb" \
  --lr_scheduler="constant" \
  --lr_warmup_steps=0 \
  --max_train_steps=500 \
  --validation_prompt="A photo of sks dog in a bucket" \
  --validation_epochs=25 \
  --seed="0" \
  --push_to_hub

To better track our training experiments, we're using the following flags in the command above:

report_to="wandb will ensure the training runs are tracked on Weights and Biases. To use it, be sure to install wandb with pip install wandb.
validation_prompt and validation_epochs to allow the script to do a few validation inference runs. This allows us to qualitatively check if the training is progressing as expected.

Our experiments were conducted on a single 40GB A100 GPU.

Dog toy example with < 16GB VRAM

By making use of gradient_checkpointing (which is natively supported in Diffusers), xformers, and bitsandbytes libraries, you can train SDXL LoRAs with less than 16GB of VRAM by adding the following flags to your accelerate launch command:

+  --enable_xformers_memory_efficient_attention \
+  --gradient_checkpointing \
+  --use_8bit_adam \
+  --mixed_precision="fp16" \

and making sure that you have the following libraries installed:

bitsandbytes>=0.40.0
xformers>=0.0.20

Inference

Once training is done, we can perform inference like so:

from huggingface_hub.repocard import RepoCard
from diffusers import DiffusionPipeline
import torch

lora_model_id = <"lora-sdxl-dreambooth-id">
card = RepoCard.load(lora_model_id)
base_model_id = card.data.to_dict()["base_model"]

pipe = DiffusionPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16)
pipe = pipe.to("cuda")
pipe.load_lora_weights(lora_model_id)
image = pipe("A picture of a sks dog in a bucket", num_inference_steps=25).images[0]
image.save("sks_dog.png")

We can further refine the outputs with the Refiner:

from huggingface_hub.repocard import RepoCard
from diffusers import DiffusionPipeline, StableDiffusionXLImg2ImgPipeline
import torch

lora_model_id = <"lora-sdxl-dreambooth-id">
card = RepoCard.load(lora_model_id)
base_model_id = card.data.to_dict()["base_model"]

# Load the base pipeline and load the LoRA parameters into it. 
pipe = DiffusionPipeline.from_pretrained(base_model_id, torch_dtype=torch.float16)
pipe = pipe.to("cuda")
pipe.load_lora_weights(lora_model_id)

# Load the refiner.
refiner = StableDiffusionXLImg2ImgPipeline.from_pretrained(
    "stabilityai/stable-diffusion-xl-refiner-1.0", torch_dtype=torch.float16, use_safetensors=True, variant="fp16"
)
refiner.to("cuda")

prompt = "A picture of a sks dog in a bucket"
generator = torch.Generator("cuda").manual_seed(0)

# Run inference.
image = pipe(prompt=prompt, output_type="latent", generator=generator).images[0]
image = refiner(prompt=prompt, image=image[None, :], generator=generator).images[0]
image.save("refined_sks_dog.png")

Here's a side-by-side comparison of the with and without Refiner pipeline outputs:

Without Refiner	With Refiner

Training with text encoder(s)

Alongside the UNet, LoRA fine-tuning of the text encoders is also supported. To do so, just specify --train_text_encoder while launching training. Please keep the following points in mind:

SDXL has two text encoders. So, we fine-tune both using LoRA.
When not fine-tuning the text encoders, we ALWAYS precompute the text embeddings to save memory.

Specifying a better VAE

SDXL's VAE is known to suffer from numerical instability issues. This is why we also expose a CLI argument namely --pretrained_vae_model_name_or_path that lets you specify the location of a better VAE (such as this one).

Notes

In our experiments, we found that SDXL yields good initial results without extensive hyperparameter tuning. For example, without fine-tuning the text encoders and without using prior-preservation, we observed decent results. We didn't explore further hyper-parameter tuning experiments, but we do encourage the community to explore this avenue further and share their results with us 🤗

Results

You can explore the results from a couple of our internal experiments by checking out this link: https://wandb.ai/sayakpaul/dreambooth-lora-sd-xl. Specifically, we used the same script with the exact same hyperparameters on the following datasets: