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license: mit |
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## Pony Diffusion XL Model Card |
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Pony Diffusion XL is a latent text-to-image diffusion model capable of generating images of Horses, mainly, and other things. For more information about how Stable Diffusion functions, please have a look at 🤗's [Stable Diffusion blog](https://huggingface.co/blog/stable_diffusion). |
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You can use this with the 🧨Diffusers library from [Hugging Face](https://huggingface.co). |
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### Diffusers |
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```py |
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from diffusers import StableDiffusionPipeline |
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import torch |
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pipeline = StableDiffusionPipeline.from_pretrained("nroggendorff/ponyxl").to("cuda") |
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image = pipeline(prompt="a chibi doll").images[0] |
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image.save("horse.png") |
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``` |
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### Model Details |
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- `train_batch_size`: 1 |
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- `gradient_accumulation_steps`: 4 |
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- `learning_rate`: 1e-2 |
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- `lr_warmup_steps`: 500 |
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- `mixed_precision`: "fp16" |
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- `eval_metric`: "mean_squared_error" |
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### Limitations |
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- The model does not achieve perfect photorealism |
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- The model cannot render legible text |
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### Developed by |
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- Noa Linden Roggendorff |
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*This model card was written by Noa Roggendorff and is based on the [Stable Diffusion v1-5 Model Card](https://huggingface.co/runwayml/stable-diffusion-v1-5).* |
