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--- |
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license: openrail++ |
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base_model: stabilityai/stable-diffusion-xl-base-1.0 |
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tags: |
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- stable-diffusion-xl |
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- stable-diffusion-xl-diffusers |
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- text-to-image |
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- diffusers |
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- inpainting |
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inference: false |
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--- |
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# SD-XL Inpainting 0.1 Model Card |
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![inpaint-example](inpaint-examples-min.png) |
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SD-XL Inpainting 0.1 is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input, with the extra capability of inpainting the pictures by using a mask. |
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The SD-XL Inpainting 0.1 was initialized with the `stable-diffusion-xl-base-1.0` weights. The model is trained for 40k steps at resolution 1024x1024 and 5% dropping of the text-conditioning to improve classifier-free classifier-free guidance sampling. For inpainting, the UNet has 5 additional input channels (4 for the encoded masked-image and 1 for the mask itself) whose weights were zero-initialized after restoring the non-inpainting checkpoint. During training, we generate synthetic masks and, in 25% mask everything. |
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## How to use |
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```py |
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from diffusers import AutoPipelineForInpainting |
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from diffusers.utils import load_image |
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import torch |
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pipe = AutoPipelineForInpainting.from_pretrained("diffusers/stable-diffusion-xl-1.0-inpainting-0.1", torch_dtype=torch.float16, variant="fp16").to("cuda") |
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img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png" |
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mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png" |
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image = load_image(img_url).resize((1024, 1024)) |
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mask_image = load_image(mask_url).resize((1024, 1024)) |
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prompt = "a tiger sitting on a park bench" |
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generator = torch.Generator(device="cuda").manual_seed(0) |
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image = pipe( |
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prompt=prompt, |
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image=image, |
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mask_image=mask_image, |
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guidance_scale=8.0, |
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num_inference_steps=20, # steps between 15 and 30 work well for us |
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strength=0.99, # make sure to use `strength` below 1.0 |
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generator=generator, |
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).images[0] |
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``` |
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**How it works:** |
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`image` | `mask_image` |
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:-------------------------:|:-------------------------:| |
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<img src="https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png" alt="drawing" width="300"/> | <img src="https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png" alt="drawing" width="300"/> |
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`prompt` | `Output` |
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:-------------------------:|:-------------------------:| |
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<span style="position: relative;bottom: 150px;">a tiger sitting on a park bench</span> | <img src="https://huggingface.co/datasets/valhalla/images/resolve/main/tiger.png" alt="drawing" width="300"/> |
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## Model Description |
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- **Developed by:** The Diffusers team |
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- **Model type:** Diffusion-based text-to-image generative model |
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- **License:** [CreativeML Open RAIL++-M License](https://huggingface.co/stabilityai/stable-diffusion-xl-base-1.0/blob/main/LICENSE.md) |
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- **Model Description:** This is a model that can be used to generate and modify images based on text prompts. It is a [Latent Diffusion Model](https://arxiv.org/abs/2112.10752) that uses two fixed, pretrained text encoders ([OpenCLIP-ViT/G](https://github.com/mlfoundations/open_clip) and [CLIP-ViT/L](https://github.com/openai/CLIP/tree/main)). |
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## Uses |
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### Direct Use |
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The model is intended for research purposes only. Possible research areas and tasks include |
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- Generation of artworks and use in design and other artistic processes. |
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- Applications in educational or creative tools. |
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- Research on generative models. |
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- Safe deployment of models which have the potential to generate harmful content. |
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- Probing and understanding the limitations and biases of generative models. |
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Excluded uses are described below. |
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### Out-of-Scope Use |
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The model was not trained to be factual or true representations of people or events, and therefore using the model to generate such content is out-of-scope for the abilities of this model. |
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## Limitations and Bias |
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### Limitations |
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- The model does not achieve perfect photorealism |
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- The model cannot render legible text |
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- The model struggles with more difficult tasks which involve compositionality, such as rendering an image corresponding to “A red cube on top of a blue sphere” |
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- Faces and people in general may not be generated properly. |
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- The autoencoding part of the model is lossy. |
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- When the strength parameter is set to 1 (i.e. starting in-painting from a fully masked image), the quality of the image is degraded. The model retains the non-masked contents of the image, but images look less sharp. We're investing this and working on the next version. |
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### Bias |
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While the capabilities of image generation models are impressive, they can also reinforce or exacerbate social biases. |
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