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---
library_name: diffusers
license: mit
---
## Pony Diffusion XL Model Card
Pony Diffusion XL is a latent text-to-image diffusion model capable of generating images of Horses, mainly, and other things. For more information about how Stable Diffusion functions, please have a look at 🤗's [Stable Diffusion blog](https://huggingface.co/blog/stable_diffusion).
You can use this with the 🧨Diffusers library from [Hugging Face](https://huggingface.co).
### Diffusers
```py
from diffusers import StableDiffusionXLPipeline
import torch
pipeline = StableDiffusionXLPipeline.from_pretrained("nroggendorff/ponyxl").to("cuda")
image = pipeline(prompt="a chibi doll").images[0]
image.save("horse.png")
```
### Model Details
- `train_batch_size`: 1
- `gradient_accumulation_steps`: 4
- `learning_rate`: 1e-2
- `lr_warmup_steps`: 500
- `mixed_precision`: "fp16"
- `eval_metric`: "mean_squared_error"
### Limitations
- The model does not achieve perfect photorealism
- The model cannot render legible text
### Developed by
- Noa Linden Roggendorff
*This model card was written by Noa Roggendorff and is based on the [Stable Diffusion v1-5 Model Card](https://huggingface.co/runwayml/stable-diffusion-v1-5).* |